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196110 | https://en.wikipedia.org/wiki/Hooded%20crow | Hooded crow | The hooded crow (Corvus cornix), also called the scald-crow or hoodie, is a Eurasian bird species in the genus Corvus. Widely distributed, it is found across Northern, Eastern, and Southeastern Europe, as well as parts of the Middle East. It is an ashy grey bird with black head, throat, wings, tail, and thigh feathers, as well as a black bill, eyes, and feet. Like other corvids, it is an omnivorous and opportunistic forager and feeder.
The hooded crow is so similar in morphology and habits to the carrion crow (Corvus corone) that for many years they were considered by most authorities to be geographical races of one species. Hybridization observed where their ranges overlapped added weight to this view. However, since 2002, the hooded crow has been elevated to full species status after closer observation; the hybridisation was less than expected and hybrids had decreased vigour. Within the hooded crow species, four subspecies are recognized.
Taxonomy
The hooded crow was one of the many species originally described by Carl Linnaeus in his landmark 1758 10th edition of Systema Naturae; he gave it the binomial name Corvus cornix. Linnaeus specified the type locality as "Europa", but this was restricted to Sweden by the German ornithologist Ernst Hartert in 1903. The genus name Corvus is Latin for "raven" while the specific epithet cornix is Latin for "crow". The hooded crow was subsequently considered a subspecies of the carrion crow for many years, hence known as Corvus corone cornix, due to similarities in structure and habits. "Hooded crow" has been designated as the official name by the International Ornithologists' Union (IOC). It is locally known as a 'hoodie craw' or simply 'hoodie' in Scotland and as a grey crow in Northern Ireland. It is also known locally as "Scotch crow" and "Danish crow". In Irish, it is called , or the "grey crow", as its name also means in the Slavic languages and in Danish. It is referred to as the "fog crow" () in German, and the "dolman crow" () in Hungarian.
Subspecies
Four subspecies of the hooded crow are now recognised; previously, all were considered subspecies of Corvus corone. A fifth subspecies, C. c. sardonius (Kleinschmidt, 1903) has been listed, although it has been alternately partitioned between C. c. sharpii (most populations), C. c. cornix (Corsican population), and the Middle Eastern C. c. pallescens.
C. c. cornix Linnaeus, 1758 – the nominate race, occurs in Britain, Ireland and the rest of Europe south to Corsica.
C. c. sharpii Oates, 1889 – named for English zoologist Richard Bowdler Sharpe. This is a paler grey form found from western Siberia through to the Caucasus region and Iran.
C. c. pallescens (Madarász, 1904) – found in Turkey and Egypt, and is a paler form as its name suggests.
C. c. capellanus Sclater, PL, 1877 – sometimes considered a separate species. This distinctive form occurs in Iraq and southwestern Iran. It has very pale grey plumage, which looks almost white from a distance. It is possibly distinct enough to be considered a separate species.
Genetic difference from carrion crows
The hooded crow (Corvus cornix) and carrion crow (Corvus corone) are two closely related species whose geographical distribution across Europe is illustrated in the accompanying diagram. It is believed that this distribution might have resulted from the glaciation cycles during the Pleistocene, which caused the parent population to split into isolates which subsequently re-expanded their ranges when the climate warmed causing secondary contact. Jelmer Poelstra and coworkers sequenced almost the entire genomes of both species in populations at varying distances from the contact zone to find that the two species were nearly genetically identical, both in their DNA and in its expression (in the form of mRNA), except for the lack of expression of a small portion (<0.28%) of the genome (situated on avian chromosome 18) in the hooded crow, which imparts the lighter plumage colouration on its torso. Thus the two species can viably hybridize, and occasionally do so at the contact zone, but the all-black carrion crows on the one side of the contact zone mate almost exclusively with other all-black carrion crows, while the same occurs among the hooded crows on the other side of the contact zone. They concluded that it was only the outward appearance of the two species that inhibits hybridization.
Description
Except for the head, throat, wings, tail, and thigh feathers, which are black and mostly glossy, the plumage of the hooded crow is ash-grey, with the dark shafts giving it a streaky appearance. The bill and legs are black; the iris dark brown. Only one moult occurs, in autumn, as in other crow species. Male hooded crows tend to be larger than females, although the two sexes are otherwise similar in appearance. Their flight is slow, heavy and usually straight. Their length varies from . When first hatched, the young are much blacker than the parents. Juveniles have duller plumage with bluish or greyish eyes, and initially possess a red mouth. Wingspan is and weight is on average 510 g.
The hooded crow, with its contrasted greys and blacks, is visually distinct from both the carrion crow and the rook, but the call notes of the hooded and carrion crows are almost indistinguishable.
Distribution
The hooded crow breeds in northern and eastern Europe, and closely allied forms inhabit southern Europe and western Asia. Where its range overlaps with that of the carrion crow, as in northern Britain, Germany, Denmark, northern Italy, and Siberia, their hybrids are fertile. However, the hybrids are less well-adapted than purebred birds (one of the reasons behind its reclassification as a distinct species from the carrion crow).
In the British Isles, the hooded crow breeds regularly in Scotland, the Isle of Man, and the Scottish Islands; it also breeds widely in Ireland. In autumn, some migratory birds arrive on the east coast of Britain. In the past, this was a more common visitor.
Behaviour
Diet
The hooded crow is omnivorous, with a diet similar to that of the carrion crow, and is a constant scavenger. It drops molluscs and crabs to break them after the manner of the carrion crow, to the point that an old Scottish name for empty sea urchin shells was "crow's cups". On coastal cliffs, the eggs of gulls, cormorants, and other birds are stolen when their owners are absent, and hooded crows will enter the burrows of puffins to steal eggs. It will also feed on small mammals, scraps, smaller birds, and carrion. The hooded crow often hides food to feed on later, especially meat, nuts, and any insects that may be present on these, in places such as rain gutters, flower pots, or in the earth under bushes. Other crows will often watch a crow that hides food and then search the hiding place later when the first crow has left.
Nesting
Nesting occurs later in colder regions: mid-May to mid-June in northwest Russia, Shetland, and the Faroe Islands, and late February in the Persian Gulf region. In warmer parts of the European Archipelago, the clutch is laid in April. The bulky stick nest is normally placed in a tall tree, but cliff ledges, old buildings, and pylons may be used. Nests are occasionally placed on or near the ground. The nest resembles that of the carrion crow, but on the coast, seaweed is often interwoven in the structure, and animal bones and wire are also frequently incorporated. The four to six brown-speckled blue eggs are in size and weigh , of which 6% is shell. The altricial young are incubated for 17–19 days by the female alone, that is fed by the male. They fledge after 32 to 36 days. Incubating females have been reported to obtain most of their own food and later that for their young.
The typical lifespan is unknown, but that of the carrion crow is four years. The maximum recorded age for a hooded crow is 16 years, and 9 months.
This species is a secondary host of the parasitic great spotted cuckoo, the European magpie being the preferred host. However, in areas where the latter species is absent, such as Israel and Egypt, the hooded crow becomes the normal corvid host.
This species, like its relative, is regularly killed by farmers and on grouse estates. In County Cork, Ireland, the county's gun clubs shot over 23,000 hooded crows in two years in the early 1980s. Since 1981, they have been protected under the Wildlife and Countryside Act, meaning it is illegal to knowingly kill, injure, or capture them.
Status
The IUCN Red List does not distinguish the hooded crow from the carrion crow, but the two species together have an extensive range, estimated at , and a large population, including an estimated 14 to 34 million individuals in Europe alone. They are not believed to approach the thresholds for the population decline criterion of the IUCN Red List (i.e., declining more than 30% in ten years or three generations), so are evaluated as least concern. The carrion crow and hooded crow hybrid zone is slowly spreading northwest, but the hooded crow has on the order of three million territories in just Europe (excluding Russia). This movement is also attested to by the fact that in April 2020 the hooded crow was redlisted in Sweden, where the Species Information Centre does distinguish between hooded and carrion crow.
Cultural significance
In Irish folklore, the bird appears on the shoulder of the dying Cú Chulainn, and could also be a manifestation of the Morrígan, the wife of Tethra, or the Cailleach. This idea has persisted, and the hooded crow is associated with fairies in the Scottish highlands and Ireland; in the 18th century, Scottish shepherds would make offerings to them to keep them from attacking sheep. In Faroese folklore, a maiden would go out on Candlemas morn and throw a stone, then a bone, then a clump of turf at a hooded crow – if it flew over the sea, her husband would be a foreigner; if it landed on a farm or house, she would marry a man from there, but if it stayed put, she would remain unmarried.
The old name of Royston crow originates from the days when this bird was a common winter visitor to southern England, with the sheep fields around Royston, Hertfordshire providing carcasses on which the birds could feed. The local newspaper, founded in 1855, is called The Royston Crow, and the hooded crow also features on the town's coat of arms.
The hooded crow is one of the 37 Norwegian birds depicted in the Bird Room of the Royal Palace in Oslo. Jethro Tull mentions the hooded crow on the song "Jack Frost and the Hooded Crow" as a bonus track on the digitally remastered version of Broadsword and the Beast and on their The Christmas Album.
In January 2014, a hooded crow and a yellow-legged gull each attacked one of two peace doves which Pope Francis had allowed children to release in Vatican City.
| Biology and health sciences | Corvoidea | Animals |
196121 | https://en.wikipedia.org/wiki/Platelet | Platelet | Platelets or thrombocytes () are a blood component whose function (along with the coagulation factors) is to react to bleeding from blood vessel injury by clumping, thereby initiating a blood clot. Platelets have no cell nucleus; they are fragments of cytoplasm derived from the megakaryocytes of the bone marrow or lung, which then enter the circulation. Platelets are found only in mammals, whereas in other vertebrates (e.g. birds, amphibians), thrombocytes circulate as intact mononuclear cells.
One major function of platelets is to contribute to hemostasis: the process of stopping bleeding at the site of interrupted endothelium. They gather at the site and, unless the interruption is physically too large, they plug the hole. First, platelets attach to substances outside the interrupted endothelium: adhesion. Second, they change shape, turn on receptors and secrete chemical messengers: activation. Third, they connect to each other through receptor bridges: aggregation. Formation of this platelet plug (primary hemostasis) is associated with activation of the coagulation cascade, with resultant fibrin deposition and linking (secondary hemostasis). These processes may overlap: the spectrum is from a predominantly platelet plug, or "white clot" to a predominantly fibrin, or "red clot" or the more typical mixture. Berridge adds retraction and platelet inhibition as fourth and fifth steps, while others would add a sixth step, wound repair. Platelets participate in both innate and adaptive intravascular immune responses.
In addition to facilitating the clotting process, platelets contain cytokines and growth factors which can promote wound healing and regeneration of damaged tissues.
Term
The term thrombocyte (clot cell) came into use in the early 1900s and is sometimes used as a synonym for platelet; but not generally in the scientific literature, except as a root word for other terms related to platelets (e.g. thrombocytopenia meaning low platelets). The term thrombocytes are proper for mononuclear cells found in the blood of non-mammalian vertebrates: they are the functional equivalent of platelets, but circulate as intact cells rather than cytoplasmic fragments of bone marrow megakaryocytes.
In some contexts, the word thrombus is used interchangeably with the word clot, regardless of its composition (white, red, or mixed). In other contexts it is used to contrast a normal from an abnormal clot: thrombus arises from physiologic hemostasis, thrombosis arises from a pathologic and excessive quantity of clot. In a third context it is used to contrast the result from the process: thrombus is the result, thrombosis is the process.
Morphology
Structure
Structurally the platelet can be divided into four zones, from peripheral to innermost:
Peripheral zone — rich in glycoproteins required for platelet adhesion, activation and aggregation. For example, GPIb/IX/V; GPVI; GPIIb/IIIa
Sol-gel zone — rich in microtubules and microfilaments, allowing platelets to maintain a discoid shape
Organelle zone — rich in platelet granules. Alpha granules contain clotting mediators such as factor V, factor VIII, fibrinogen, fibronectin, platelet-derived growth factor, and chemotactic agents. Delta granules, or dense bodies, contain ADP, calcium, and serotonin, which are platelet-activating mediators.
Membranous zone — membranes derived from megakaryocyte smooth endoplasmic reticulum organized into a dense tubular system that is responsible for thromboxane A2 synthesis. This dense tubular system is connected to the surface platelet membrane to aid thromboxane A2 release.
Shape
Circulating inactivated platelets are biconvex discoid (lens-shaped) structures, 2–3 μm in greatest diameter. Activated platelets have cell membrane projections covering their surface.
In a first approximation, the shape can be considered similar to oblate spheroids, with a semiaxis ratio of 2 to 8. This approximation can be used to model the hydrodynamic and optical properties of a population, as well as to restore the geometric parameters of individual measured platelets by flow cytometry. More accurate biophysical models of platelet surface morphology that model its shape from first principles, make it possible to obtain a more realistic platelet geometry in a calm and activated state.
Development
Megakaryocyte and platelet production is regulated by thrombopoietin, a hormone produced in the kidneys and liver.
Each megakaryocyte produces between 1,000 and 3,000 platelets during its lifetime.
An average of 1011 platelets are produced daily in a healthy adult.
Reserve platelets are stored in the spleen and are released when needed by splenic contraction induced by the sympathetic nervous system.
The average life span of circulating platelets is 8 to 9 days. Life span of individual platelets is controlled by the internal apoptotic regulating pathway, which has a Bcl-xL timer.
Old platelets are destroyed by phagocytosis in the spleen and liver.
Hemostasis
The fundamental function of platelets is to clump together to stop acute bleeding. This process is complex, as more than 193 proteins and 301 interactions are involved in platelet dynamics. Despite much overlap, platelet function can be modeled in three steps:
Adhesion
Thrombus formation on an intact endothelium is prevented by nitric oxide, prostacyclin, and CD39.
Endothelial cells attach to the subendothelial collagen by von Willebrand factor (VWF), which these cells produce. VWF is also stored in the Weibel-Palade bodies of the endothelial cells and secreted constitutively into the blood. Platelets store vWF in their alpha granules.
When the endothelial layer is disrupted, collagen and VWF anchor platelets to the subendothelium. Platelet GP1b-IX-V receptor binds with VWF; and GPVI receptor and integrin α2β1 bind with collagen.
Activation
Inhibition
The intact endothelial lining inhibits platelet activation by producing nitric oxide, endothelial-ADPase, and PGI2 (prostacyclin). Endothelial-ADPase degrades the platelet activator ADP.
Resting platelets maintain active calcium efflux via a cyclic AMP-activated calcium pump. Intracellular calcium concentration determines platelet activation status, as it is the second messenger that drives platelet conformational change and degranulation. Endothelial prostacyclin binds to prostanoid receptors on the surface of resting platelets. This event stimulates the coupled Gs protein to increase adenylate cyclase activity and increases the production of cAMP, further promoting the efflux of calcium and reducing intracellular calcium availability for platelet activation.
ADP on the other hand binds to purinergic receptors on the platelet surface. Since the thrombocytic purinergic receptor P2Y12 is coupled to Gi proteins, ADP reduces platelet adenylate cyclase activity and cAMP production, leading to accumulation of calcium inside the platelet by inactivating the cAMP calcium efflux pump. The other ADP-receptor P2Y1 couples to Gq that activates phospholipase C-beta 2 (PLCB2), resulting in inositol 1,4,5-trisphosphate (IP3) generation and intracellular release of more calcium. This together induces platelet activation. Endothelial ADPase degrades ADP and prevents this from happening. Clopidogrel and related antiplatelet medications also work as purinergic receptor P2Y12 antagonists. Data suggest that ADP activates the PI3K/Akt pathway during a first wave of aggregation, leading to thrombin generation and PAR‐1 activation, which evokes a second wave of aggregation.
Trigger (induction)
Platelet activation begins seconds after adhesion occurs. It is triggered when collagen from the subendothelium binds with its receptors (GPVI receptor and integrin α2β1) on the platelet. GPVI is associated with the Fc receptor gamma chain and leads via the activation of a tyrosine kinase cascade finally to the activation of PLC-gamma2 (PLCG2) and more calcium release.
Tissue factor also binds to factor VII in the blood, which initiates the extrinsic coagulation cascade to increase thrombin production. Thrombin is a potent platelet activator, acting through Gq and G12. These are G protein-coupled receptors and they turn on calcium-mediated signaling pathways within the platelet, overcoming the baseline calcium efflux. Families of three G proteins (Gq, Gi, G12) operate together for full activation. Thrombin also promotes secondary fibrin-reinforcement of the platelet plug. Platelet activation in turn degranulates and releases factor V and fibrinogen, potentiating the coagulation cascade. Platelet plugging and coagulation occur simultaneously, with each inducing the other to form the final fibrin-crosslinked thrombus.
Components (consequences)
GPIIb/IIIa activation
Collagen-mediated GPVI signalling increases the platelet production of thromboxane A2 (TXA2) and decreases the production of prostacyclin. This occurs by altering the metabolic flux of platelet's eicosanoid synthesis pathway, which involves enzymes phospholipase A2, cyclo-oxygenase 1, and thromboxane-A synthase. Platelets secrete thromboxane A2, which acts on the platelet's own thromboxane receptors on the platelet surface (hence the so-called "out-in" mechanism), and those of other platelets. These receptors trigger intraplatelet signaling, which converts GPIIb/IIIa receptors to their active form to initiate aggregation.
Granule secretion
Platelets contain dense granules, lambda granules, and alpha granules. Activated platelets secrete the contents of these granules through their canalicular systems to the exterior. Bound and activated platelets degranulate to release platelet chemotactic agents to attract more platelets to the site of endothelial injury. Granule characteristics:
α granules (alpha granules) — containing P-selectin, platelet factor 4, transforming growth factor-β1, platelet-derived growth factor, fibronectin, B-thromboglobulin, vWF, fibrinogen, and coagulation factors V and XIII
δ granules (delta or dense granules) — containing ADP or ATP, calcium, and serotonin
γ granules (gamma granules) — similar to lysosomes and contain several hydrolytic enzymes
λ granules (lambda granules) — contents involved in resorption during later stages of vessel repair
Morphology change
As shown by flow cytometry and electron microscopy, the most sensitive sign of activation, when exposed to platelets using ADP, are morphological changes. Mitochondrial hyperpolarization is a key event in initiating morphology changes. Intraplatelet calcium concentration increases, stimulating the interplay between the microtubule/actin filament complex. The continuous changes in shape from the unactivated to the fully activated platelet are best seen via scanning electron microscopy. The three steps along this path are named early dendritic, early spread, and spread. The surface of the unactivated platelet looks similar to the surface of the brain–a wrinkled appearance from numerous shallow folds that increase the surface area; early dendritic, an octopus with multiple arms and legs; early spread, an uncooked frying egg in a pan, the "yolk" is the central body; and the spread, a cooked fried egg with a denser central body.
These changes are all brought about by the interaction of the microtubule/actin complex with the platelet cell membrane and open canalicular system (OCS), which is an extension and invagination of that membrane. This complex runs just beneath these membranes and is the chemical motor that pulls the invaginated OCS out of the interior of the platelet, like turning pants pockets inside out, creating the dendrites. This process is similar to the mechanism of contraction in a muscle cell. The entire OCS thus becomes indistinguishable from the initial platelet membrane as it forms the "fried egg". This dramatic increase in surface area comes about with neither stretching nor adding phospholipids to the platelet membrane.
Platelet-coagulation factor interactions: coagulation facilitation
Platelet activation causes its membrane surface to become negatively charged. One of the signaling pathways turns on scramblase, which moves negatively charged phospholipids from the inner to the outer platelet membrane surface. These phospholipids then bind the tenase and prothrombinase complexes, two of the sites of interplay between platelets and the coagulation cascade. Calcium ions are essential for the binding of these coagulation factors.
In addition to interacting with vWF and fibrin, platelets interact with thrombin, Factors X, Va, VIIa, XI, IX, and prothrombin to complete formation via the coagulation cascade. Human platelets do not express tissue factor. Rat platelets do express tissue factor protein and carry both tissue factor pre-mRNA and mature mRNA.
Aggregation
Platelet aggregation begins minutes after activation, and occurs as a result of turning on the GPIIb/IIIa receptor, allowing these receptors to bind with vWF or fibrinogen. Each platelet has around 60,000 of these receptors. When any one or more of at least nine different platelet surface receptors are turned on during activation, intraplatelet signaling pathways cause existing GpIIb/IIIa receptors to change shape — curled to straight — and thus become capable of binding.
Since fibrinogen is a rod-like protein with nodules on either end capable of binding GPIIb/IIIa, activated platelets with exposed GPIIb/IIIa can bind fibrinogen to aggregate. GPIIb/IIIa may also further anchor the platelets to subendothelial vWF for additional structural stabilisation.
Classically it was thought that this was the only mechanism involved in aggregation, but three other mechanisms have been identified which can initiate aggregation, depending on the velocity of blood flow (i.e. shear range).
Immune function
Platelets have a central role in innate immunity, initiating and participating in multiple inflammatory processes, directly binding and even destroying pathogens. Clinical data show that many patients with serious bacterial or viral infections have thrombocytopenia, thus reducing their contribution to inflammation. Platelet-leukocyte aggregates (PLAs) found in circulation are typical in sepsis or inflammatory bowel disease, showing the connection between thrombocytes and immune cells.
The platelet cell membrane has receptors for collagen. Following rupture of the blood vessel wall, platelets are exposed and adhere to the collagen in the surrounding tissue.
Immunothrombosis
As hemostasis is a basic function of thrombocytes in mammals, it also has its uses in possible infection confinement. In case of injury, platelets, together with the coagulation cascade, provide the first line of defense by forming a blood clot. Hemostasis and host defense were thus intertwined in evolution. For example, in the Atlantic horseshoe crab (estimated to be over 400 million years old), the only blood cell type, the amebocyte, facilitates both the hemostatic function and the encapsulation and phagocytosis of pathogens by means of exocytosis of intracellular granules containing bactericidal defense molecules. Blood clotting supports immune function by trapping the bacteria.
Although thrombosis, blood coagulation in intact blood vessels, is usually viewed as a pathological immune response, leading to obturation of lumen of blood vessel and subsequent hypoxic tissue damage, in some cases, directed thrombosis, called immunothrombosis, can locally control the spread of an infection. The thrombosis is directed in concordance of platelets, neutrophils and monocytes. The process is initiated either by immune cells by activating their pattern recognition receptors (PRRs), or by platelet-bacterial binding. Platelets can bind to bacteria either directly through thrombocytic PRRs and bacterial surface proteins, or via plasma proteins that bind both to platelets and bacteria. Monocytes respond to bacterial pathogen-associated molecular patterns (PAMPs), or damage-associated molecular patterns (DAMPs) by activating the extrinsic pathway of coagulation. Neutrophils facilitate the blood coagulation by NETosis, while platelets facilitate neutrophils' NETosis. NETs bind tissue factor, binding the coagulation centers to the location of infection. They also activate the intrinsic coagulation pathway by providing its negatively charged surface to the factor XII. Other neutrophil secretions, such as proteolytic enzymes which cleave coagulation inhibitors, also bolster the process.
In case of imbalance throughout the regulation of immunothrombosis, this process can become aberrant. Regulatory defects in immunothrombosis are suspected to be a major factor in pathological thrombosis in forms such as disseminated intravascular coagulation (DIC) or deep vein thrombosis. DIC in sepsis is a prime example of both the dysregulated coagulation process as well as an undue systemic inflammatory response, resulting in a multitude of microthrombi of similar composition to that in physiological immunothrombosis — fibrin, platelets, neutrophils and NETs.
Inflammation
Platelets rapidly deploy to sites of injury or infection, and potentially modulate inflammatory processes by interacting with leukocytes and secreting cytokines, chemokines, and other inflammatory mediators. Platelets also secrete platelet-derived growth factor (PDGF).
Platelets modulate neutrophils by forming platelet-leukocyte aggregates (PLAs). These formations induce upregulated production of αmβ2 (Mac-1) integrin in neutrophils. Interaction with PLAs also induces degranulation and increased phagocytosis in neutrophils.
Platelets are the largest source of soluble CD40L which induces production of reactive oxygen species (ROS) and upregulate expression of adhesion molecules, such as E-selectin, ICAM-1, and VCAM-1, in neutrophils, activates macrophages and activates cytotoxic response in T and B lymphocytes.
Mammalian platelets lacking nucleus are able to conduct autonomous locomotion. Platelets are active scavengers, scaling walls of blood vessels and reorganising the thrombus. They are able to recognize and adhere to many surfaces, including bacteria, and can envelop them in their open canalicular system (OCP), leading to a proposal to name the process as covercytosis (OCS) rather than phagocytosis, as OCS is merely an invagination of outer plasma membrane. These platelet-bacteria bundles provide an interaction platform for neutrophils that destroy bacteria using the NETosis and phagocytosis.
Platelets also participate in chronic inflammatory disease, such as synovitis or rheumatoid arthritis. Platelets are activated by collagen receptor glycoprotein IV (GPVI). Proinflammatory platelet microvesicles trigger constant cytokine secretion from neighboring fibroblast-like synoviocytes, most prominently Il-6 and Il-8. Inflammatory damage to the surrounding extracellular matrix continuously reveals more collagen, maintaining microvesicle production.
Adaptive immunity
Activated platelets are able to participate in adaptive immunity, interacting with antibodies. They are able to specifically bind IgG through FcγRIIA, a receptor for IgG's constant fragment (Fc). When activated and bound to IgG opsonised bacteria, platelets release reactive oxygen species (ROS), antimicrobial peptides, defensins, kinocidins and proteases, killing the bacteria directly. Platelets also secrete proinflammatory and procoagulant mediators such as inorganic polyphosphates or platelet factor 4 (PF4), connecting innate and adaptive immune responses.
Signs and symptoms of disorders
Spontaneous and excessive bleeding can occur because of platelet disorders. This bleeding can be caused by deficient numbers of platelets, dysfunctional platelets, or platelet densities over 1 million/microliter. (The excessive numbers create a relative von Willebrand factor deficiency due to sequestration.)
Bleeding due to a platelet disorder or a coagulation factor disorder can be distinguished by the characteristics and location of the bleeding. Platelet bleeding involves bleeding from a cut that is prompt and excessive, but can be controlled by pressure; spontaneous bleeding into the skin which causes a purplish stain named by its size: petechiae, purpura, ecchymoses; bleeding into mucous membranes causing bleeding gums, nose bleed, and gastrointestinal bleeding; menorrhagia; and intraretinal and intracranial bleeding.
Excessive numbers of platelets, and/or normal platelets responding to abnormal vessel walls, can result in venous thrombosis and arterial thrombosis. The symptoms depend on the thrombosis site.
Measurement and testing
Measurement
Platelet concentration in the blood (i.e. platelet count), can be measured manually using a hemocytometer, or by placing blood in an automated platelet analyzer using particle counting, such as a Coulter counter or optical methods. Most common blood testing methods include platelet count in their measurements, usually reported as PLT.
Platelet concentrations vary between individuals and over time, with the population average between 250,000 and 260,000 cells per mm3 (equivalent to per microliter), but the typical laboratory accepted normal range is between 150,000 and 400,000 cells per mm3 or 150–400 × 109 per liter.
On a stained blood smear, platelets appear as dark purple spots, about 20% of the diameter of red blood cells. The smear reveals size, shape, qualitative number, and clumping. A healthy adult typically has 10 to 20 times more red blood cells than platelets.
Bleeding time
Bleeding time was developed as a test of platelet function by Duke in 1910. Duke's test measured the time taken for bleeding to stop from a standardized wound in the ear lobe that was blotted every 30 seconds, considering less than 3 minutes as normal. Bleeding time has low sensitivity and specificity for mild to moderate platelet disorders and is no longer recommended for screening.
Multiple electrode aggregometry
In multiple electrode aggregometry, anticoagulated whole blood is mixed with saline and a platelet agonist in a single-use cuvette with two pairs of electrodes. The increase in impedance between the electrodes as platelets aggregate onto them, is measured and visualized as a curve.
Light transmission aggregometry
In light transmission aggregometry (LTA), platelet-rich plasma is placed between a light source and a photocell. Unaggregated plasma allows relatively little light to pass through. After adding an agonist, the platelets aggregate, increasing light transmission, which is detected by a photocell.
Whole blood impedance aggregometry
Whole blood impedance aggregometry (WBA) measures the change in electrical impedance between two electrodes when platelet aggregation is induced by an agonist. Whole blood lumiaggregometry may increase the test sensitivity to impairment of platelet granule secretion.
PFA-100
The PFA-100 (Platelet Function Assay — 100) is a system for analysing platelet function in which citrated whole blood is aspirated through a disposable cartridge containing an aperture within a membrane coated with either collagen and epinephrine or collagen and ADP. These agonists induce platelet adhesion, activation and aggregation, leading to rapid occlusion of the aperture and cessation of blood flow termed the closure time (CT). An elevated CT with EPI and collagen can indicate intrinsic defects such as von Willebrand disease, uremia, or circulating platelet inhibitors. A follow-up test involving collagen and ADP is used to indicate if the abnormal CT with collagen and EPI was caused by the effects of acetyl sulfosalicylic acid (aspirin) or medications containing inhibitors. The PFA-100 is highly sensitive to von Willebrand disease, but is only moderately sensitive to defects in platelet function.
Disorders
Low platelet concentration is called thrombocytopenia, and is due to either decreased production or increased destruction. Elevated platelet concentration is called thrombocytosis, and is either congenital, reactive (to cytokines), or due to unregulated production: one of the myeloproliferative neoplasms or certain other myeloid neoplasms. A disorder of platelet function is called a thrombocytopathy or a platelet function disorder.
Normal platelets can respond to an abnormality on the vessel wall rather than to hemorrhage, resulting in inappropriate platelet adhesion/activation and thrombosis: the formation of a clot within an intact vessel. This type of thrombosis arises by mechanisms different from those of a normal clot: extending the fibrin of venous thrombosis; extending an unstable or ruptured arterial plaque, causing arterial thrombosis; and microcirculatory thrombosis. An arterial thrombus may partially obstruct blood flow, causing downstream ischemia, or may completely obstruct it, causing downstream tissue death.:
The three broad categories of platelet disorders are "not enough", "dysfunctional", and "too many".
Thrombocytopenia
Immune thrombocytopenia (ITP) — formerly known as immune thrombocytopenic purpura and idiopathic thrombocytopenic purpura
Splenomegaly
Gaucher's disease
Familial thrombocytopenia
Chemotherapy
Babesiosis
Dengue fever
Onyalai
Thrombotic thrombocytopenic purpura
HELLP syndrome
Hemolytic–uremic syndrome
Drug-induced thrombocytopenic purpura (five known drugs — most problematic is heparin-induced thrombocytopenia (HIT)
Pregnancy-associated
Neonatal alloimmune associated
Aplastic anemia
Transfusion-associated
Pseudothrombocytopenia
Vaccine-induced immune thrombotic thrombocytopenia (VITT)
Altered platelet function (thrombocytopathy)
Congenital
Disorders of adhesion
Bernard–Soulier syndrome
Disorders of activation
Disorders of granule amount or release
Hermansky–Pudlak syndrome
Gray platelet syndrome
ADP receptor defect
Decreased cyclooxygenase activity
Platelet storage pool deficiency
Disorders of aggregation
Glanzmann's thrombasthenia
Wiskott–Aldrich syndrome
Disorders of coagulant activity
COAT platelet defect
Scott syndrome
Acquired
Disorders of adhesion
Paroxysmal nocturnal hemoglobinuria
Asthma
Aspirin-exacerbated respiratory disease (AERD/Samter's triad)
Cancer
Malaria
Decreased cyclooxygenase activity
Thrombocytosis and thrombocythemia
Reactive
Chronic infection
Chronic inflammation
Malignancy
Hyposplenism (post-splenectomy)
Iron deficiency
Acute blood loss
Myeloproliferative neoplasms — platelets are both elevated and activated
Essential thrombocythemia
Polycythemia vera
Associated with other myeloid neoplasms
Congenital
Pharmacology
Anti-inflammatory drugs
Some drugs used to treat inflammation have the unwanted side effect of suppressing normal platelet function. These are the non-steroidal anti-inflammatory drugs (NSAIDS). Aspirin irreversibly disrupts platelet function by inhibiting cyclooxygenase-1 (COX1), and hence normal hemostasis. The resulting platelets are unable to produce new cyclooxygenase because they have no DNA. Normal platelet function does not return until the use of aspirin has ceased and enough of the affected platelets have been replaced by new ones, which can take over a week. Ibuprofen, another NSAID, does not have such a long duration effect, with platelet function usually returning within 24 hours, and taking ibuprofen before aspirin prevents the irreversible effects of aspirin.
Drugs that suppress platelet function
These drugs are used to prevent thrombus formation.
Oral agents
Aspirin
Cilostazol
Clopidogrel
Prasugrel
Ticagrelor
Ticlopidine
Drugs that stimulate platelet production
Desmopressin
Factor VIIa
Thrombopoietin mimetics
Intravenous agents
Abciximab
Eptifibatide
Tirofiban
Others: oprelvekin, romiplostim, eltrombopag, argatroban
Therapies
Transfusion
Indications
Platelet transfusion is most frequently used to correct unusually low platelet counts, either to prevent spontaneous bleeding (typically at counts below 10×109/L) or in anticipation of medical procedures that necessarily involve some bleeding. For example, in patients undergoing surgery, a level below 50×109/L is associated with abnormal surgical bleeding, and regional anaesthetic procedures such as epidurals are avoided for levels below 80×109/L. Platelets may also be transfused when the platelet count is normal but the platelets are dysfunctional, such as when an individual is taking aspirin or clopidogrel. Finally, platelets may be transfused as part of a massive transfusion protocol, in which the three major blood components (red blood cells, plasma, and platelets) are transfused to address severe hemorrhage. Platelet transfusion is contraindicated in thrombotic thrombocytopenic purpura (TTP), as it fuels the coagulopathy. Platelet transfusion is generally ineffective, and thus contraindicated, for prophylaxis in immune thrombocytopenia (ITP), because the transfused platelets are immediately cleared; however, it is indicated to treat bleeding.
Collection
Platelets are either isolated from collected units of whole blood and pooled to make a therapeutic dose, or collected by platelet apheresis: blood is taken from the donor, passed through a device which removes the platelets, and the remainder is returned to the donor in a closed loop. The industry standard is for platelets to be tested for bacteria before transfusion to avoid septic reactions, which can be fatal. Recently the AABB Industry Standards for Blood Banks and Transfusion Services (5.1.5.1) has allowed use of pathogen reduction technology as an alternative to bacterial screenings in platelets.
Pooled whole-blood platelets, sometimes called "random" platelets, are separated by one of two methods. In the US, a unit of whole blood is placed into a large centrifuge in what is referred to as a "soft spin". At these settings, the platelets remain suspended in the plasma. The platelet-rich plasma (PRP) is removed from the red cells, then centrifuged at a faster setting to harvest the platelets from the plasma. In other regions of the world, the unit of whole blood is centrifuged using settings that cause the platelets to become suspended in the "buffy coat" layer, which includes the platelets and the white blood cells. The "buffy coat" is isolated in a sterile bag, suspended in a small amount of red blood cells and plasma, then centrifuged again to separate the platelets and plasma from the red and white blood cells. Regardless of the initial method of preparation, multiple donations may be combined into one container using a sterile connection device to manufacture a single product with the desired therapeutic dose.
Apheresis platelets are collected using a mechanical device that draws blood from the donor and centrifuges the collected blood to separate out the platelets and other components to be collected. The remaining blood is returned to the donor. The advantage to this method is that a single donation provides at least one therapeutic dose, as opposed to the multiple donations for whole-blood platelets. This means that a recipient is exposed to fewer donors and has less risk of transfusion-transmitted disease and other complications. Sometimes a person such as a cancer patient who requires routine transfusions of platelets receives repeated donations from a specific donor to minimize risk. Pathogen reduction of platelets using for example, riboflavin and UV light treatments can reduce the infectious load of pathogens contained in donated blood products. Another photochemical treatment process utilizing amotosalen and UVA light has been developed for the inactivation of viruses, bacteria, parasites, and leukocytes. In addition, apheresis platelets tend to contain fewer contaminating red blood cells because the collection method is more efficient than "soft spin" centrifugation.
Storage
Platelets collected by either method have a typical shelf life of five days. This results in supply shortages, as testing donations often requires up to a full day. No effective preservative solutions have been devised for platelets.
Platelets are stored under constant agitation at . Units cannot be refrigerated as this causes platelets to change shape and lose function. Storage at room temperature provides an environment where any introduced bacteria may proliferate and subsequently cause bacteremia. The United States requires products to be tested for the presence of bacterial contamination before transfusion.
Delivery
Platelets do not need to belong to the same A-B-O blood group as the recipient or be cross-matched to ensure immune compatibility between donor and recipient unless they contain a significant amount of red blood cells (RBCs). The presence of RBCs imparts a reddish-orange color to the product and is usually associated with whole-blood platelets. Some sites may type platelets, but this is not critical.
Prior to issuing platelets to the recipient, they may be irradiated to prevent transfusion-associated graft versus host disease or they may be washed to remove the plasma.
The change in the recipient's platelet count after transfusion is termed the "increment" and is calculated by subtracting the pre-transfusion platelet count from the post-transfusion count. Many factors affect the increment including body size, the number of platelets transfused, and clinical features that may cause premature destruction of the transfused platelets. When recipients fail to demonstrate an adequate post-transfusion increment, this is termed platelet transfusion refractoriness.
Platelets, either apheresis-derived or random-donor, can be processed through a volume reduction process. In this process, the platelets are spun in a centrifuge and plasma is removed, leaving 10 to 100 mL of platelet concentrate. Such volume-reduced platelets are normally transfused only to neonatal and pediatric patients when a large volume of plasma could overload the child's small circulatory system. The lower volume of plasma also reduces the chances of an adverse transfusion reaction to plasma proteins. Volume reduced platelets have a shelf life of four hours.
Wound repair
The blood clot is only a temporary solution to stop bleeding; tissue repair is needed. Small interruptions in the endothelium are handled by physiological mechanisms; large interruptions by a trauma surgeon. The fibrin is slowly dissolved by the fibrinolytic enzyme, plasmin, and the platelets are cleared by phagocytosis.
Platelets release platelet-derived growth factor (PDGF), a potent chemotactic agent; and TGF beta, which stimulates the deposition of extracellular matrix; fibroblast growth factor, insulin-like growth factor 1, platelet-derived epidermal growth factor, and vascular endothelial growth factor. Local application of these factors in increased concentrations through platelet-rich plasma (PRP) is used as an adjunct in wound healing.
Non-mammals
Instead of platelets, non-mammalian vertebrates have nucleated thrombocytes, which resemble B lymphocytes in morphology. They aggregate in response to thrombin, but not to ADP, serotonin, nor adrenaline, as platelets do.
History
George Gulliver in 1841 drew pictures of platelets using the twin lens (compound) microscope invented in 1830 by Joseph Jackson Lister. This microscope improved resolution sufficiently to make it possible to see platelets for the first time.
William Addison in 1842 drew pictures of a platelet-fibrin clot.
Lionel Beale in 1864 was the first to publish a drawing showing platelets.
Max Schultze in 1865 described what he called "spherules", which he noted were much smaller than red blood cells, occasionally clumped, and were sometimes found in collections of fibrin material.
Giulio Bizzozero in 1882 studied the blood of amphibians microscopically in vivo. He named Schultze's spherules (It.) piastrine: little plates. Bizzozero possibly proposed the name Blutplattchen.
William Osler observed platelets and, in published lectures in 1886, called them a third corpuscle and a blood plaque; and described them as "a colorless protoplasmic disc".
James Wright examined blood smears using the stain named for him, and used the term plates in his 1906 publication, changing to platelets in his 1910 publication.
| Biology and health sciences | Circulatory system | Biology |
196130 | https://en.wikipedia.org/wiki/Bone%20marrow | Bone marrow | Bone marrow is a semi-solid tissue found within the spongy (also known as cancellous) portions of bones. In birds and mammals, bone marrow is the primary site of new blood cell production (or haematopoiesis). It is composed of hematopoietic cells, marrow adipose tissue, and supportive stromal cells. In adult humans, bone marrow is primarily located in the ribs, vertebrae, sternum, and bones of the pelvis. Bone marrow comprises approximately 5% of total body mass in healthy adult humans, such that a man weighing 73 kg (161 lbs) will have around 3.7 kg (8 lbs) of bone marrow.
Human marrow produces approximately 500 billion blood cells per day, which join the systemic circulation via permeable vasculature sinusoids within the medullary cavity. All types of hematopoietic cells, including both myeloid and lymphoid lineages, are created in bone marrow; however, lymphoid cells must migrate to other lymphoid organs (e.g. thymus) in order to complete maturation.
Bone marrow transplants can be conducted to treat severe diseases of the bone marrow, including certain forms of cancer such as leukemia. Several types of stem cells are related to bone marrow. Hematopoietic stem cells in the bone marrow can give rise to hematopoietic lineage cells, and mesenchymal stem cells, which can be isolated from the primary culture of bone marrow stroma, can give rise to bone, adipose, and cartilage tissue.
Structure
The composition of marrow is dynamic, as the mixture of cellular and non-cellular components (connective tissue) shifts with age and in response to systemic factors. In humans, marrow is colloquially characterized as "red" or "yellow" marrow (, , respectively) depending on the prevalence of hematopoietic cells vs fat cells. While the precise mechanisms underlying marrow regulation are not understood, compositional changes occur according to stereotypical patterns. For example, a newborn baby's bones exclusively contain hematopoietically active "red" marrow, and there is a progressive conversion towards "yellow" marrow with age. In adults, red marrow is found mainly in the central skeleton, such as the pelvis, sternum, cranium, ribs, vertebrae and scapulae, and variably found in the proximal epiphyseal ends of long bones such as the femur and humerus. In circumstances of chronic hypoxia, the body can convert yellow marrow back to red marrow to increase blood cell production.
Hematopoietic components
At the cellular level, the main functional component of bone marrow includes the progenitor cells which are destined to mature into blood and lymphoid cells. Human marrow produces approximately 500 billion blood cells per day. Marrow contains hematopoietic stem cells which give rise to the three classes of blood cells that are found in circulation: white blood cells (leukocytes), red blood cells (erythrocytes), and platelets (thrombocytes).
Stroma
The stroma of the bone marrow includes all tissue not directly involved in the marrow's primary function of hematopoiesis. Stromal cells may be indirectly involved in hematopoiesis, providing a microenvironment that influences the function and differentiation of hematopoietic cells. For instance, they generate colony stimulating factors, which have a significant effect on hematopoiesis. Cell types that constitute the bone marrow stroma include:
fibroblasts (reticular connective tissue)
macrophages, which contribute especially to red blood cell production, as they deliver iron for hemoglobin production.
adipocytes (fat cells)
osteoblasts (synthesize bone)
osteoclasts (resorb bone)
endothelial cells, which form the sinusoids. These derive from endothelial stem cells, which are also present in the bone marrow.
Function
Central hematopoietic and antigen-responsive organ
That bone marrow is a priming site for T-cell responses to blood-borne antigens was first described in 2003. Mature circulating naïve T cells home to bone marrow sinuses after they have passed through arteries and arterioles. They transmigrate sinus endothelium and enter the parenchyma which contains dendritic cells (DCs). These have a capacity of antigen uptake, processing, and presentation. Cognate interactions between antigen-specific T cells and antigen-presenting DCs (APCs) in parenchyma lead to rapid T-APC cluster formation followed by T cell activation, T cell proliferation and T cell re-circulation to blood. These findings were corroborated and extended in 2013 by in situ two-photon dynamic imaging of mice skulls.
Importance for storage and long-term survival of memory B and memory T cells
Bone marrow is a nest for migratory memory T cells and a sanctuary for plasma cells. This has implications for adaptive immunity and vaccinology. Memory B and T cells persist in the parenchyma in dedicated survival niches organized by stromal cells. This memory can be maintained over long time periods in the form of quiescent cells or by repeated antigenic restimulation. Bone marrow protects and optimizes immunological memory during dietary restriction. In cancer patients, cancer-reactive memory T cells can arise in bone marrow spontaneously or after specific vaccination. Bone marrow is a center of a variety of immune activities: i) hematopoiesis, ii) osteogenesis, iii) immune responses, iv) distinction between self and non-self antigens, v) central immune regulatory function, vi) storage of memory cells, vii) immune surveillance of the central nervous system, viii) adaptation to energy crisis, ix) provision of mesenchymal stem cells for tissue repair.
Mesenchymal stem cells
The bone marrow stroma contains mesenchymal stem cells (MSCs), which are also known as marrow stromal cells. These are multipotent stem cells that can differentiate into a variety of cell types. MSCs have been shown to differentiate, in vitro or in vivo, into osteoblasts, chondrocytes, myocytes, marrow adipocytes and beta-pancreatic islets cells.
Bone marrow barrier
The blood vessels of the bone marrow constitute a barrier, inhibiting immature blood cells from leaving the marrow. Only mature blood cells contain the membrane proteins, such as aquaporin and glycophorin, that are required to attach to and pass the blood vessel endothelium. Hematopoietic stem cells may also cross the bone marrow barrier, and may thus be harvested from blood.
Lymphatic role
The red bone marrow is a key element of the lymphatic system, being one of the primary lymphoid organs that generate lymphocytes from immature hematopoietic progenitor cells. The bone marrow and thymus constitute the primary lymphoid tissues involved in the production and early selection of lymphocytes. Furthermore, bone marrow performs a valve-like function to prevent the backflow of lymphatic fluid in the lymphatic system.
Compartmentalization
Biological compartmentalization is evident within the bone marrow, in that certain cell types tend to aggregate in specific areas. For instance, erythrocytes, macrophages, and their precursors tend to gather around blood vessels, while granulocytes gather at the borders of the bone marrow.
As food
People have used animal bone-marrow in cuisine worldwide for millennia, as in the famed Milanese Ossobuco.
Clinical significance
Disease
The normal bone marrow architecture can be damaged or displaced by aplastic anemia, malignancies such as multiple myeloma, or infections such as tuberculosis, leading to a decrease in the production of blood cells and blood platelets. The bone marrow can also be affected by various forms of leukemia, which attacks its hematologic progenitor cells. Furthermore, exposure to radiation or chemotherapy will kill many of the rapidly dividing cells of the bone marrow, and will therefore result in a depressed immune system. Many of the symptoms of radiation poisoning are due to damage sustained by the bone marrow cells.
To diagnose diseases involving the bone marrow, a bone marrow aspiration is sometimes performed. This typically involves using a hollow needle to acquire a sample of red bone marrow from the crest of the ilium under general or local anesthesia.
Imaging
Medical imaging may provide a limited amount of information regarding bone marrow. Plain film x-rays pass through soft tissues such as marrow and do not provide visualization, although any changes in the structure of the associated bone may be detected. CT imaging has somewhat better capacity for assessing the marrow cavity of bones, although with low sensitivity and specificity. For example, normal fatty "yellow" marrow in adult long bones is of low density (-30 to -100 Hounsfield units), between subcutaneous fat and soft tissue. Tissue with increased cellular composition, such as normal "red" marrow or cancer cells within the medullary cavity will measure variably higher in density.
MRI is more sensitive and specific for assessing bone composition. MRI enables assessment of the average molecular composition of soft tissues and thus provides information regarding the relative fat content of marrow. In adult humans, "yellow" fatty marrow is the dominant tissue in bones, particularly in the (peripheral) appendicular skeleton. Because fat molecules have a high T1-relaxivity, T1-weighted imaging sequences show "yellow" fatty marrow as bright (hyperintense). Furthermore, normal fatty marrow loses signal on fat-saturation sequences, in a similar pattern to subcutaneous fat.
When "yellow" fatty marrow becomes replaced by tissue with more cellular composition, this change is apparent as decreased brightness on T1-weighted sequences. Both normal "red" marrow and pathologic marrow lesions (such as cancer) are darker than "yellow" marrow on T1-weight sequences, although can often be distinguished by comparison with the MR signal intensity of adjacent soft tissues. Normal "red" marrow is typically equivalent or brighter than skeletal muscle or intervertebral disc on T1-weighted sequences.
Fatty marrow change, the inverse of red marrow hyperplasia, can occur with normal aging, though it can also be seen with certain treatments such as radiation therapy. Diffuse marrow T1 hypointensity without contrast enhancement or cortical discontinuity suggests red marrow conversion or myelofibrosis. Falsely normal marrow on T1 can be seen with diffuse multiple myeloma or leukemic infiltration when the water to fat ratio is not sufficiently altered, as may be seen with lower grade tumors or earlier in the disease process.
Histology
Bone marrow examination is the pathologic analysis of samples of bone marrow obtained via biopsy and bone marrow aspiration. Bone marrow examination is used in the diagnosis of a number of conditions, including leukemia, multiple myeloma, anemia, and pancytopenia. The bone marrow produces the cellular elements of the blood, including platelets, red blood cells and white blood cells. While much information can be gleaned by testing the blood itself (drawn from a vein by phlebotomy), it is sometimes necessary to examine the source of the blood cells in the bone marrow to obtain more information on hematopoiesis; this is the role of bone marrow aspiration and biopsy.
The ratio between myeloid series and erythroid cells is relevant to bone marrow function, and also to diseases of the bone marrow and peripheral blood, such as leukemia and anemia. The normal myeloid-to-erythroid ratio is around 3:1; this ratio may increase in myelogenous leukemias, decrease in polycythemias, and reverse in cases of thalassemia.
Donation and transplantation
In a bone marrow transplant, hematopoietic stem cells are removed from a person and infused into another person (allogenic) or into the same person at a later time (autologous). If the donor and recipient are compatible, these infused cells will then travel to the bone marrow and initiate blood cell production. Transplantation from one person to another is conducted for the treatment of severe bone marrow diseases, such as congenital defects, autoimmune diseases or malignancies. The patient's own marrow is first killed off with drugs or radiation, and then the new stem cells are introduced. Before radiation therapy or chemotherapy in cases of cancer, some of the patient's hematopoietic stem cells are sometimes harvested and later infused back when the therapy is finished to restore the immune system.
Bone marrow stem cells can be induced to become neural cells to treat neurological illnesses, and can also potentially be used for the treatment of other illnesses, such as inflammatory bowel disease. In 2013, following a clinical trial, scientists proposed that bone marrow transplantation could be used to treat HIV in conjunction with antiretroviral drugs; however, it was later found that HIV remained in the bodies of the test subjects.
Harvesting
The stem cells are typically harvested directly from the red marrow in the iliac crest, often under general anesthesia. The procedure is minimally invasive and does not require stitches afterwards. Depending on the donor's health and reaction to the procedure, the actual harvesting can be an outpatient procedure, or can require 1–2 days of recovery in the hospital.
Another option is to administer certain drugs that stimulate the release of stem cells from the bone marrow into circulating blood. An intravenous catheter is inserted into the donor's arm, and the stem cells are then filtered out of the blood. This procedure is similar to that used in blood or platelet donation. In adults, bone marrow may also be taken from the sternum, while the tibia is often used when taking samples from infants. In newborns, stem cells may be retrieved from the umbilical cord.
Fertility Aid
One of the most damaged areas of the body following chemotherapy is typically the uterus. Following uterine damage due to cancer treatment, follicle damage makes it difficult for individuals to get pregnant even if viable ova are present. With bone marrow stem cell transplants, chemotherapy patients have been able to increase their fertility as follicle damage is repaired. As follicles are necessary for ovum attachment to the endometrium, it is important for these areas to be repaired in order to increase fertility. For individuals who have sustained egg and follicle damage, IVF has been found to be more effective following bone marrow stem cell transplantation.
One human clinical case has shown improvements of uterine lining thickness and overall endometrium repair following bone marrow stem cell transplantation. This repair allowed for the patient to successfully become pregnant and carry to term.
Additional repairs following bone marrow stem cell transplant to the endometrium include increased vascularity and iron levels, with egg implantation clustering around areas with high blood flow.
Persistent viruses
Using quantitative polymerase chain reaction (qPCR) and next-generation sequencing (NGS) a maximum of five DNA viruses per individual have been identified. Included were several herpesviruses, hepatitis B virus, Merkel cell polyomavirus, and human papillomavirus 31. Given the reactivation and/or oncogenic potential of these viruses, their repercussion on hematopoietic and malignant disorders calls for further studies.
Fossil record
The earliest fossilised evidence of bone marrow was discovered in 2014 in Eusthenopteron, a lobe-finned fish which lived during the Devonian period approximately 370 million years ago. Scientists from Uppsala University and the European Synchrotron Radiation Facility used X-ray synchrotron microtomography to study the fossilised interior of the skeleton's humerus, finding organised tubular structures akin to modern vertebrate bone marrow. Eusthenopteron is closely related to the early tetrapods, which ultimately evolved into the land-dwelling mammals and lizards of the present day.
| Biology and health sciences | Skeletal system | Biology |
196147 | https://en.wikipedia.org/wiki/House%20sparrow | House sparrow | The house sparrow (Passer domesticus) is a bird of the sparrow family Passeridae, found in most parts of the world. It is a small bird that has a typical length of and a mass of . Females and young birds are coloured pale brown and grey, and males have brighter black, white, and brown markings. One of about 25 species in the genus Passer, the house sparrow is native to most of Europe, the Mediterranean Basin, and a large part of Asia. Its intentional or accidental introductions to many regions, including parts of Australasia, Africa, and the Americas, make it the most widely distributed wild bird.
The house sparrow is strongly associated with human habitation, and can live in urban or rural settings. Though found in widely varied habitats and climates, it typically avoids extensive woodlands, grasslands, polar regions, and hot, dry deserts far away from human development. For sustenance, the house sparrow routinely feeds at home and public bird feeding stations, but naturally feeds on the seeds of grains, flowering plants and weeds. However, it is an opportunistic, omnivorous eater, and commonly catches insects, their larvae, caterpillars, invertebrates and many other natural foods.
Because of its numbers, ubiquity, and association with human settlements, the house sparrow is culturally prominent. It is extensively, and usually unsuccessfully, persecuted as an agricultural pest. It has also often been kept as a pet, as well as being a food item and a symbol of lust, sexual potency, commonness, and vulgarity. Though it is widespread and abundant, its numbers have declined in some areas. The animal's conservation status is listed as least concern on the IUCN Red List.
Description
Measurements and shape
The house sparrow is typically about long, ranging from . The house sparrow is a compact bird with a full chest and a large, rounded head. Its bill is stout and conical with a culmen length of , strongly built as an adaptation for eating seeds. Its tail is short, at long. The wing chord is , and the tarsus is . Wingspan ranges from .
In mass, the house sparrow ranges from . Females usually are slightly smaller than males. The median mass on the European continent for both sexes is about , and in more southerly subspecies is around . Younger birds are smaller, males are larger during the winter, and females are larger during the breeding season. Birds at higher latitudes, colder climates, and sometimes higher altitudes are larger (under Bergmann's rule), both between and within subspecies.
Plumage
The plumage of the house sparrow is mostly different shades of grey and brown. The sexes exhibit strong dimorphism: the female is mostly buffish above and below, while the male has boldly coloured head markings, a reddish back, and grey underparts. The male has a dark grey crown from the top of its bill to its back, and chestnut brown flanking its crown on the sides of its head. It has black around its bill, on its throat, and on the spaces between its bill and eyes (lores). It has a small white stripe between the lores and crown and small white spots immediately behind the eyes (postoculars), with black patches below and above them. The underparts are pale grey or white, as are the cheeks, ear coverts, and stripes at the base of the head. The upper back and mantle are a warm brown, with broad black streaks, while the lower back, rump and upper tail coverts are greyish brown.
The male is duller in fresh nonbreeding plumage, with whitish tips on many feathers. Wear and preening expose many of the bright brown and black markings, including most of the black throat and chest patch, called the "bib" or "badge". The badge is variable in width and general size, and may signal social status or fitness. This hypothesis has led to a "veritable 'cottage industry of studies, which have only conclusively shown that patches increase in size with age. The male's bill is dark grey, but black in the breeding season.
The female has no black markings or grey crown. Its upperparts and head are brown with darker streaks around the mantle and a distinct pale supercilium. Its underparts are pale grey-brown. The female's bill is brownish-grey and becomes darker in breeding plumage approaching the black of the male's bill.
Juveniles are similar to the adult female, but deeper brown below and paler above, with paler and less defined supercilia. Juveniles have broader buff feather edges, and tend to have looser, scruffier plumage, like moulting adults. Juvenile males tend to have darker throats and white postoculars like adult males, while juvenile females tend to have white throats. However, juveniles cannot be reliably sexed by plumage: some juvenile males lack any markings of the adult male, and some juvenile females have male features. The bills of young birds are light yellow to straw, paler than the female's bill. Immature males have paler versions of the adult male's markings, which can be very indistinct in fresh plumage. By their first breeding season, young birds generally are indistinguishable from other adults, though they may still be paler during their first year.
Voice
Most house sparrow vocalisations are variations on its short and frequent chirping call. Transcribed as chirrup, tschilp, or philip, this note is made as a contact call by flocking or resting birds; or by males to proclaim nest ownership and invite pairing. In the breeding season, the male gives this call repetitively, with emphasis and speed, but not much rhythm, forming what is described either as a song or an "ecstatic call" similar to a song. Young birds also give a true song, especially in captivity, a warbling similar to that of the European greenfinch.
Aggressive males give a trilled version of their call, transcribed as "chur-chur-r-r-it-it-it-it". This call is also used by females in the breeding season, to establish dominance over males while displacing them to feed young or incubate eggs. House sparrows give a nasal alarm call, the basic sound of which is transcribed as quer, and a shrill chree call in great distress. Another vocalisation is the "appeasement call", a soft quee given to inhibit aggression, usually given between birds of a mated pair. These vocalisations are not unique to the house sparrow, but are shared, with small variations, by all sparrows.
Variation
Some variation is seen in the 12 subspecies of house sparrows, which are divided into two groups, the Oriental P. d. indicus group, and the Palaearctic P. d. domesticus group. Birds of the P. d. domesticus group have grey cheeks, while P. d. indicus group birds have white cheeks, as well as bright colouration on the crown, a smaller bill, and a longer black bib. The subspecies P. d. tingitanus differs little from the nominate subspecies, except in the worn breeding plumage of the male, in which the head is speckled with black and underparts are paler. P. d. balearoibericus is slightly paler than the nominate, but darker than P. d. bibilicus. P. d. bibilicus is paler than most subspecies, but has the grey cheeks of P. d. domesticus group birds. The similar P. d. persicus is paler and smaller, and P. d. niloticus is nearly identical but smaller. Of the less widespread P. d. indicus group subspecies, P. d. hyrcanus is larger than P. d. indicus, P. d. hufufae is paler, P. d. bactrianus is larger and paler, and P. d. parkini is larger and darker with more black on the breast than any other subspecies.
Identification
The house sparrow can be confused with a number of other seed-eating birds, especially its relatives in the genus Passer. Many of these relatives are smaller, with an appearance that is neater or "cuter", as with the Dead Sea sparrow. The light brown-coloured female can often not be distinguished from other females, and is nearly identical to those of the Spanish and Italian sparrows. The Eurasian tree sparrow is smaller and slenderer with a chestnut crown and a black patch on each cheek. The male Spanish sparrow and Italian sparrow are distinguished by their chestnut crowns. The Sind sparrow is very similar but smaller, with less black on the male's throat and a distinct pale supercilium on the female.
Taxonomy and systematics
Names
The house sparrow was among the first animals to be given a scientific name in the modern system of biological classification, since it was described by Carl Linnaeus, in the 1758 10th edition of Systema Naturae. It was described from a type specimen collected in Sweden, with the name Fringilla domestica. Later, the genus name Fringilla came to be used only for the common chaffinch and its relatives, and the house sparrow has usually been placed in the genus Passer created by French zoologist Mathurin Jacques Brisson in 1760.
The bird's scientific name and its usual English name have the same meaning. The Latin word , like the English word "sparrow", is a term for small active birds, coming from a root word referring to speed. The Latin word domesticus means "belonging to the house", like the common name a reference to its association with humans. The house sparrow is also called by a number of alternative English names, including English sparrow, chiefly in North America; and Indian sparrow or Indian house sparrow, for the birds of the Indian subcontinent and Central Asia. Dialectal names include sparr, sparrer, spadger, spadgick, and philip, mainly in southern England; spug and spuggy, mainly in northern England; spur and sprig, mainly in Scotland; and spatzie or spotsie, from the German , in North America.
Taxonomy
The genus Passer contains about 25 species, depending on the authority, 26 according to the Handbook of the Birds of the World. Most Passer species are dull-coloured birds with short, square tails and stubby, conical beaks, between long. Mitochondrial DNA studies suggest that speciation in the genus occurred during the Pleistocene and earlier, while other evidence suggests speciation occurred 25,000 to 15,000 years ago. Within Passer, the house sparrow is part of the "Palaearctic black-bibbed sparrows" group and a close relative of the Mediterranean "willow sparrows".
The taxonomy of the house sparrow and its Mediterranean relatives is complicated. The common type of "willow sparrow" is the Spanish sparrow, which resembles the house sparrow in many respects. It frequently prefers wetter habitats than the house sparrow, and it is often colonial and nomadic. In most of the Mediterranean, one or both species occur, with some degree of hybridisation. In North Africa, the two species hybridise extensively, forming highly variable mixed populations with a full range of characters from pure house sparrows to pure Spanish sparrows.
In most of Italy, the breeding species is the Italian sparrow, which has an appearance intermediate between those of the house and Spanish sparrows. Its specific status and origin are the subject of much debate, but it may be a case of long-ago hybrid speciation. In the Alps, the Italian sparrow intergrades over a narrow roughly strip with the house sparrow, and some house sparrows migrate into the Italian sparrow's range in winter. On the Mediterranean islands of Malta, Gozo, Crete, Rhodes, and Karpathos, other apparently intermediate birds are of unknown status.
Subspecies
A large number of subspecies have been named, of which 12 were recognised in the Handbook of the Birds of the World. These subspecies are divided into two groups, the Palaearctic P. d. domesticus group, and the Oriental P. d. indicus group. Several Middle Eastern subspecies, including P. d. biblicus, are sometimes considered a third, intermediate group. The subspecies P. d. indicus was described as a species, and was considered to be distinct by many ornithologists during the 19th century.
Migratory birds of the subspecies P. d. bactrianus in the P. d. indicus group were recorded overlapping with P. d. domesticus birds without hybridising in the 1970s, so the Soviet scientists Edward I. Gavrilov and M. N. Korelov proposed the separation of the P. d. indicus group as a separate species. However, P. d. indicus group and P. d. domesticus group birds intergrade in a large part of Iran, so this split is rarely recognised.
In North America, house sparrow populations are more differentiated than those in Europe. This variation follows predictable patterns, with birds at higher latitudes being larger and darker and those in arid areas being smaller and paler. However, how much this is caused by evolution or by environment is not clear. Similar observations have been made in New Zealand and in South Africa. The introduced house sparrow populations may be distinct enough to merit subspecies status, especially in North America and southern Africa, and American ornithologist Harry Church Oberholser even gave the subspecies name P. d. plecticus to the paler birds of western North America.
P. d. domesticus group
P. d. domesticus Linnaeus, 1758, the nominate subspecies, is found in most of Europe, across northern Asia to Sakhalin and Kamchatka. It is the most widely introduced subspecies.
P. d. balearoibericus von Jordans, 1923, described from Majorca, is found in the Balearic Islands, southern France, the Balkans, and Anatolia.
P. d. tingitanus (Loche, 1867), described from Algeria, is found in the Maghreb from Ajdabiya in Libya to Béni Abbès in Algeria, and to Morocco's Atlantic coast. It hybridises extensively with the Spanish sparrow, especially in the eastern part of its range.
P. d. niloticus Nicoll and Bonhote, 1909, described from Faiyum, Egypt, is found along the Nile north of Wadi Halfa, Sudan. It intergrades with bibilicus in the Sinai, and with rufidorsalis in a narrow zone around Wadi Halfa. It has been recorded in Somaliland.
P. d. persicus Zarudny and Kudashev, 1916, described from the Karun River in Khuzestan, Iran, is found in the western and central Iran south of the Alborz mountains, intergrading with indicus in eastern Iran, and Afghanistan.
P. d. biblicus Hartert, 1910, described from Palestine, is found in the Middle East from Cyprus and southeastern Turkey to the Sinai in the west and from Azerbaijan to Kuwait in the east.
P. d. indicus group
P. d. hyrcanus Zarudny and Kudashev, 1916, described from Gorgan, Iran, is found along the southern coast of the Caspian Sea from Gorgan to southeastern Azerbaijan. It intergrades with P. d. persicus in the Alborz mountains, and with P. d. bibilicus to the west. It is the subspecies with the smallest range.
P. d. bactrianus Zarudny and Kudashev, 1916, described from Tashkent, is found in southern Kazakhstan to the Tian Shan and northern Iran and Afghanistan. It intergrades with persicus in Baluchistan and with indicus across central Afghanistan. Unlike most other house sparrow subspecies, it is almost entirely migratory, wintering in the plains of the northern Indian subcontinent. It is found in open country rather than in settlements, which are occupied by the Eurasian tree sparrow in its range. There is an exceptional record from Sudan.
P. d. parkini Whistler, 1920, described from Srinagar, Kashmir, is found in the western Himalayas from the Pamir Mountains to southeastern Nepal. It is migratory, like P. d. bactrianus.
P. d. indicus Jardine and Selby, 1831, described from Bangalore, is found in the Indian subcontinent south of the Himalayas, in Sri Lanka, western Southeast Asia, eastern Iran, southwestern Arabia and southern Israel.
P. d. hufufae Ticehurst and Cheeseman, 1924, described from Hofuf in Saudi Arabia, is found in northeastern Arabia.
P. d. rufidorsalis C. L. Brehm, 1855, described from Khartoum, Sudan, is found in the Nile valley from Wadi Halfa south to Renk in northern South Sudan, and in eastern Sudan, northern Ethiopia to the Red Sea coast in Eritrea. It has also been introduced to Mohéli in the Comoros.
Distribution and habitat
The house sparrow originated in the Middle East and spread, along with agriculture, to most of Eurasia and parts of North Africa. Since the mid-19th century, it has reached most of the world, chiefly due to deliberate introductions, but also through natural and shipborne dispersal. Its introduced range encompasses most of North America (including Bermuda), Central America, southern South America, southern Africa, part of West Africa, Australia, New Zealand, and islands throughout the world. It has greatly extended its range in northern Eurasia since the 1850s, and continues to do so, as was shown by its colonisation around 1990 of Iceland and Rishiri Island, Japan. The extent of its range makes it the most widely distributed wild bird on the planet.
Introduction
The house sparrow has become highly successful in most parts of the world where it has been introduced. This is mostly due to its early adaptation to living with humans, and its adaptability to a wide range of conditions. Other factors may include its robust immune response, compared to the Eurasian tree sparrow. Where introduced, it can extend its range quickly, sometimes at a rate over per year. In many parts of the world, it has been characterised as a pest, and poses a threat to native birds. A few introductions have died out or been of limited success, such as those to Greenland and Cape Verde.
The first of many successful introductions to North America occurred when birds from England were released in New York City, in 1852, intended to control the ravages of the linden moth. In North America, the house sparrow now occurs from the Northwest Territories of Canada to southern Panama, and it is one of the most abundant birds of the continent. The house sparrow was first introduced to Australia in 1863 at Melbourne and is common throughout the eastern part of the continent as far north as Cape York, but has been prevented from establishing itself in Western Australia, where every house sparrow found in the state is killed. House sparrows were introduced in New Zealand in 1859, and from there reached many of the Pacific islands, including Hawaii.
In southern Africa, birds of both the European subspecies (P. d. domesticus) and the Indian subspecies (P. d. indicus) were introduced around 1900. Birds of P. d. domesticus ancestry are confined to a few towns, while P. d. indicus birds have spread rapidly, reaching Tanzania in the 1980s. Despite this rapid spread, native relatives such as the Cape sparrow also occur and thrive in urban habitats. In South America, it was first introduced near Buenos Aires around 1870, and quickly became common in most of the southern part of the continent. It now occurs almost continuously from Tierra del Fuego to the fringes of the Amazon basin, with isolated populations as far north as coastal Venezuela.
Habitat
The house sparrow is closely associated with human habitation and cultivation. It is not an obligate commensal of humans as some have suggested: birds of the migratory Central Asian subspecies usually breed away from humans in open country, and birds elsewhere are occasionally found away from humans. The only terrestrial habitats that the house sparrow does not inhabit are dense forest and tundra. Well adapted to living around humans, it frequently lives and even breeds indoors, especially in factories, warehouses, and zoos. It has been recorded breeding in an English coal mine below ground, and feeding on the Empire State Building's observation deck at night. It reaches its greatest densities in urban centres, but its reproductive success is greater in suburbs, where insects are more abundant. On a larger scale, it is most abundant in wheat-growing areas such as the Midwestern United States.
It tolerates a variety of climates, but prefers drier conditions, especially in moist tropical climates. It has several adaptations to dry areas, including a high salt tolerance and an ability to survive without water by ingesting berries. In most of eastern Asia, the house sparrow is entirely absent, replaced by the Eurasian tree sparrow. Where these two species overlap, the house sparrow is usually more common than the Eurasian tree sparrow, but one species may replace the other in a manner that ornithologist Maud Doria Haviland described as "random, or even capricious". In most of its range, the house sparrow is extremely common, despite some declines, but in marginal habitats such as rainforest or mountain ranges, its distribution can be spotty.
Behaviour
Social behaviour
The house sparrow is a very social bird. It is gregarious during all seasons when feeding, often forming flocks with other species of birds. It roosts communally while breeding nests are usually grouped together in clumps. House sparrows also engage in social activities such as dust or water bathing and "social singing", in which birds call together in bushes. The house sparrow feeds mostly on the ground, but it flocks in trees and bushes. At feeding stations and nests, female house sparrows are dominant despite their smaller size, and they can fight over males in the breeding season.
Sleep and roosting
House sparrows sleep with the bill tucked underneath the scapular feathers. Outside of the reproductive season, they often roost communally in trees or shrubs. Much communal chirping occurs before and after the birds settle in the roost in the evening, as well as before the birds leave the roost in the morning. Some congregating sites separate from the roost may be visited by the birds prior to settling in for the night.
Body maintenance
Dust or water bathing is common and often occurs in groups. Anting is rare. Head scratching is done with the leg over the drooped wing.
Feeding
As an adult, the house sparrow mostly feeds on the seeds of grains and weeds, but it is opportunistic and adaptable, and eats whatever foods are available. In towns and cities, it often scavenges for food in garbage containers and congregates in the outdoors of restaurants and other eating establishments to feed on leftover food and crumbs. It can perform complex tasks to obtain food, such as opening automatic doors to enter supermarkets, clinging to hotel walls to watch vacationers on their balconies, and nectar robbing kowhai flowers. In common with many other birds, the house sparrow requires grit to digest the harder items in its diet. Grit can be either stone, often grains of masonry, or the shells of eggs or snails; oblong and rough grains are preferred.
Several studies of the house sparrow in temperate agricultural areas have found the proportion of seeds in its diet to be about 90%. It will eat almost any seeds, but where it has a choice, it prefers corn: oats, wheat or maize. Rural birds tend to eat more waste seed from animal dung and seed from fields while urban birds tend to eat more commercial bird seed and weed seed. In urban areas, the house sparrow also feeds largely on food provided directly or indirectly by humans, such as bread, though it prefers raw seeds. The house sparrow also eats some plant matter besides seeds, including buds, berries, and fruits such as grapes and cherries. In temperate areas, the house sparrow has an unusual habit of tearing flowers, especially yellow ones, in the spring.
Animals form another important part of the house sparrow's diet, chiefly insects, of which beetles, caterpillars, dipteran flies, and aphids are especially important. Various noninsect arthropods are eaten, as are molluscs and crustaceans where available, earthworms, and even vertebrates such as lizards and frogs. Young house sparrows are fed mostly on insects until about 15 days after hatching. They are also given small quantities of seeds, spiders, and grit. In most places, grasshoppers and crickets are the most abundant foods of nestlings. True bugs, ants, sawflies, and beetles are also important, but house sparrows take advantage of whatever foods are abundant to feed their young. House sparrows have been observed stealing prey from other birds, including American robins.
The gut microbiota of house sparrows differs between chicks and adults, with Pseudomonadota (formerly Proteobacteria) decreasing in chicks when they get to around 9 days old, whilst the relative abundance of Bacillota increase.
Locomotion
The house sparrow's flight is direct (not undulating) and flapping, averaging and about 15 wingbeats per second. On the ground, the house sparrow typically hops rather than walks. It can swim when pressed to do so by pursuit from predators. Captive birds have been recorded diving and swimming short distances under water.
Dispersal and migration
Most house sparrows do not move more than a few kilometres during their lifetimes. However, limited migration occurs in all regions. Some young birds disperse long distances, especially on coasts, and mountain birds move to lower elevations in winter. Two subspecies, P. d. bactrianus and P. d. parkini, are predominantly migratory. Unlike the birds in sedentary populations that migrate, birds of migratory subspecies prepare for migration by putting on weight.
Breeding
House sparrows can breed in the breeding season immediately following their hatching, and sometimes attempt to do so. Some birds breeding for the first time in tropical areas are only a few months old and still have juvenile plumage. Birds breeding for the first time are rarely successful in raising young, and reproductive success increases with age, as older birds breed earlier in the breeding season, and fledge more young. As the breeding season approaches, hormone releases trigger enormous increases in the size of the sexual organs and changes in day length lead males to start calling by nesting sites. The timing of mating and egg-laying varies geographically, and between specific locations and years because a sufficient supply of insects is needed for egg formation and feeding nestlings.
Males take up nesting sites before the breeding season, by frequently calling beside them. Unmated males start nest construction and call particularly frequently to attract females. When a female approaches a male during this period, the male displays by moving up and down while drooping and shivering his wings, pushing up his head, raising and spreading his tail, and showing his bib. Males may try to mate with females while calling or displaying. In response, a female will adopt a threatening posture and attack a male before flying away, pursued by the male. The male displays in front of her, attracting other males, which also pursue and display to the female. This group display usually does not immediately result in copulations. Other males usually do not copulate with the female. Copulation is typically initiated by the female giving a soft dee-dee-dee call to the male. Birds of a pair copulate frequently until the female is laying eggs, and the male mounts the female repeatedly each time a pair mates.
The house sparrow is monogamous, and typically mates for life, but birds from pairs often engage in extra-pair copulations, so about 15% of house sparrow fledglings are unrelated to their mother's mate. Males guard their mates carefully to avoid being cuckolded, and most extra-pair copulation occurs away from nest sites. Males may sometimes have multiple mates, and bigamy is mostly limited by aggression between females. Many birds do not find a nest and a mate, and instead may serve as helpers around the nest for mated pairs, a role which increases the chances of being chosen to replace a lost mate. Lost mates of both sexes can be replaced quickly during the breeding season. The formation of a pair and the bond between the two birds is tied to the holding of a nest site, though paired house sparrows can recognise each other away from the nest.
House sparrows in natural small populations, as can occur on islands, exhibit inbreeding depression. Inbreeding depression is manifested as lower survival probability and production of fewer offspring, and can occur as a result of the expression of deleterious recessive alleles. However sparrows in such populations do not appear to avoid inbreeding.
Nesting
Nest sites are varied, though cavities are preferred. Nests are most frequently built in the eaves and other crevices of houses. Holes in cliffs and banks, and tree hollows, are also used. A sparrow sometimes excavates its own nests in sandy banks or rotten branches, but more frequently uses the nests of other birds such as those of swallows in banks and cliffs, and old tree cavity nests. It usually uses deserted nests, though sometimes it usurps active ones by driving away or killing the occupants. Tree hollows are more commonly used in North America than in Europe, putting the sparrows in competition with bluebirds and other North American cavity nesters, and thereby contributing to their population declines.
Especially in warmer areas, the house sparrow may build its nests in the open, on the branches of trees, especially evergreens and hawthorns, or in the nests of large birds such as storks or magpies. In open nesting sites, breeding success tends to be lower, since breeding begins late and the nest can easily be destroyed or damaged by storms. Less common nesting sites include street lights and neon signs, favoured for their warmth; and the old open-topped nests of other songbirds, which are then domed over.
Usually the couples repeat copulation many times. Every copulation is followed by some break of 3 to 4 seconds, and in that time both pair change their position by some distance.
The nest is usually domed, though it may lack a roof in enclosed sites. It has an outer layer of stems and roots, a middle layer of dead grass and leaves, and a lining of feathers, as well as of paper and other soft materials. Nests typically have external dimensions of 20 × 30 cm (8 × 12 in), but their size varies greatly. The building of the nest is initiated by the unmated male while displaying to females. The female assists in building, but is less active than the male. Some nest building occurs throughout the year, especially after moult in autumn. In colder areas house sparrows build specially created roost nests, or roost in street lights, to avoid losing heat during the winter. House sparrows do not hold territories, but they defend their nests aggressively against intruders of the same sex.
House sparrows' nests support a wide range of scavenging insects, including nest flies such as Neottiophilum praestum, Protocalliphora blowflies, and over 1,400 species of beetle.
Eggs and young
Clutches usually comprise four or five eggs, though numbers from one to 10 have been recorded. At least two clutches are usually laid, and up to seven a year may be laid in the tropics or four a year in temperate latitudes. When fewer clutches are laid in a year, especially at higher latitudes, the number of eggs per clutch is greater. Central Asian house sparrows, which migrate and have only one clutch a year, average 6.5 eggs in a clutch. Clutch size is also affected by environmental and seasonal conditions, female age, and breeding density.
Some intraspecific brood parasitism occurs, and instances of unusually large numbers of eggs in a nest may be the result of females laying eggs in the nests of their neighbours. Such foreign eggs are sometimes recognised and ejected by females. The house sparrow is a victim of interspecific brood parasites, but only rarely, since it usually uses nests in holes too small for parasites to enter, and it feeds its young foods unsuitable for young parasites. In turn, the house sparrow has once been recorded as a brood parasite of the American cliff swallow.
The eggs are white, bluish white, or greenish white, spotted with brown or grey. Subelliptical in shape, they range from in length and in width, have an average mass of , and an average surface area of . Eggs from the tropical subspecies are distinctly smaller. Eggs begin to develop with the deposition of yolk in the ovary a few days before ovulation. In the day between ovulation and laying, egg white forms, followed by eggshell. Eggs laid later in a clutch are larger, as are those laid by larger females, and egg size is hereditary. Eggs decrease slightly in size from laying to hatching. The yolk comprises 25% of the egg, the egg white 68%, and the shell 7%. Eggs are watery, being 79% liquid, and otherwise mostly protein.
The female develops a brood patch of bare skin and plays the main part in incubating the eggs. The male helps, but can only cover the eggs rather than truly incubate them. The female spends the night incubating during this period, while the male roosts near the nest. Eggs hatch at the same time, after a short incubation period lasting 11–14 days, and exceptionally for as many as 17 or as few as 9. The length of the incubation period decreases as ambient temperature increases later in the breeding season.
Young house sparrows remain in the nest for 11 to 23 days, normally 14 to 16 days. During this time, they are fed by both parents. As newly hatched house sparrows do not have sufficient insulation, they are brooded for a few days, or longer in cold conditions. The parents swallow the droppings produced by the hatchlings during the first few days; later, the droppings are moved up to away from the nest.
The chicks' eyes open after about 4 days and, at an age of about 8 days, the young birds get their first down. If both parents perish, the ensuing intensive begging sounds of the young often attract replacement parents which feed them until they can sustain themselves. All the young in the nest leave it during the same period of a few hours. At this stage, they are normally able to fly. They start feeding themselves partly after 1 or 2 days, and sustain themselves completely after 7 to 10 days, 14 at the latest.
House sparrows in natural small populations, as can occur on islands, exhibit inbreeding depression. Inbreeding depression is manifested as lower survival probability and production of fewer offspring, and can occur as a result of the expression of deleterious recessive alleles. However sparrows in such populations do not appear to avoid inbreeding.
Survival
In adult house sparrows, annual survival is 45–65%. After fledging and leaving the care of their parents, young sparrows have a high mortality rate, which lessens as they grow older and more experienced. Only about 20–25% of birds hatched survive to their first breeding season. The oldest known wild house sparrow lived for nearly two decades; it was found dead 19 years and 9 months after it was ringed in Denmark. The oldest recorded captive house sparrow lived for 23 years. The typical ratio of males to females in a population is uncertain due to problems in collecting data, but a very slight preponderance of males at all ages is usual.
Predation
The house sparrow's main predators are cats and birds of prey, but many other animals prey on them, including corvids, squirrels, and even humans—the house sparrow has been consumed in the past by people in many parts of the world, and it still is in parts of the Mediterranean. Most species of birds of prey have been recorded preying on the house sparrow in places where records are extensive. Accipiters and the merlin in particular are major predators, though cats are likely to have a greater impact on house sparrow populations. The house sparrow is also a common victim of roadkill; on European roads, it is the bird most frequently found dead.
Parasites and disease
The house sparrow is host to a huge number of parasites and diseases, and the effect of most is unknown. Ornithologist Ted R. Anderson listed thousands, noting that his list was incomplete. The commonly recorded bacterial pathogens of the house sparrow are often those common in humans, and include Salmonella and Escherichia coli. Salmonella is common in the house sparrow, and a comprehensive study of house sparrow disease found it in 13% of sparrows tested. Salmonella epidemics in the spring and winter can kill large numbers of sparrows. The house sparrow hosts avian pox and avian malaria, which it has spread to the native forest birds of Hawaii. Many of the diseases hosted by the house sparrow are also present in humans and domestic animals, for which the house sparrow acts as a reservoir host. Arboviruses such as the West Nile virus, which most commonly infect insects and mammals, survive winters in temperate areas by going dormant in birds such as the house sparrow. A few records indicate disease extirpating house sparrow populations, especially from Scottish islands, but this seems to be rare. House sparrows are also infected by haemosporidian parasites, but less so in urban than in rural areas Toxoplasma gondii has been detected in sparrows in northwestern China where they pose a risk due to their meat being consumed in the region.
The house sparrow is infested by a number of external parasites, which usually cause little harm to adult sparrows. In Europe, the most common mite found on sparrows is Proctophyllodes, the most common ticks are Argas reflexus and Ixodes arboricola, and the most common flea on the house sparrow is Ceratophyllus gallinae. Dermanyssus blood-feeding mites are also common ectoparasites of house sparrows, and these mites can enter human habitation and bite humans, causing a condition known as gamasoidosis. A number of chewing lice occupy different niches on the house sparrow's body. Menacanthus lice occur across the house sparrow's body, where they feed on blood and feathers, while Brueelia lice feed on feathers and Philopterus fringillae occurs on the head.
Physiology
House sparrows express strong circadian rhythms of activity in the laboratory. They were among the first bird species to be seriously studied in terms of their circadian activity and photoperiodism, in part because of their availability and adaptability in captivity, but also because they can "find their way" and remain rhythmic in constant darkness.
Relationships with humans
The house sparrow is closely associated with humans. They are believed to have become associated with humans around 10,000 years ago. The Turkestan subspecies (P. d. bactrianus) is least associated with humans and considered to be evolutionarily closer to the ancestral noncommensal populations. Usually, the house sparrow is regarded as a pest, since it consumes agricultural products and spreads disease to humans and their domestic animals. Even birdwatchers often hold it in little regard because of its molestation of other birds. In most of the world, the house sparrow is not protected by law. Attempts to control house sparrows include the trapping, poisoning, or shooting of adults; the destruction of their nests and eggs; or less directly, blocking nest holes and scaring off sparrows with noise, glue, or porcupine wire. However, the house sparrow can be beneficial to humans, as well, especially by eating insect pests, and attempts at the large-scale control of the house sparrow have failed.
The house sparrow has long been used as a food item. From around 1560 to at least the 19th century in northern Europe, earthenware "sparrow pots" were hung from eaves to attract nesting birds so the young could be readily harvested. Wild birds were trapped in nets in large numbers, and sparrow pie was a traditional dish, thought, because of the association of sparrows with lechery, to have aphrodisiac properties. A traditional Indian medicine, Ciṭṭukkuruvi lēkiyam in Tamil, was sold with similar aphrodisiac claims. Sparrows were also trapped as food for falconers' birds and zoo animals. During the 1870s, there were debates on the damaging effects of sparrows in the House of Commons in England. In the early part of the 20th century, sparrow clubs culled many millions of birds and eggs in an attempt to control numbers of this perceived pest, but with only a localised impact on numbers. Sparrows were also persecuted in Germany from at least 1650 until 1970. House sparrows have been kept as pets at many times in history, though they have no bright plumage or attractive songs, and raising them is difficult. The house sparrow has an extremely large range and population, so it is assessed as least concern for conservation on the IUCN Red List.
Population decline
The IUCN estimates for the global population runs up to nearly 1.4 billion individuals, second among all wild birds perhaps only to the red-billed quelea in abundance (although the quelea is, unlike the sparrow, restricted to a single continent and has never been subject to human introductions). However, populations have been declining in many parts of the world, especially near its Eurasian places of origin. These declines were first noticed in North America, where they were initially attributed to the spread of the house finch, but have been most severe in Western Europe. Declines have even occurred in Australia, where the house sparrow was introduced recently. While no serious declines had been reported from Eastern Europe by 2006, as of 2023, a 50% decrease of the sparrow population has been registered in Bulgaria.
In Great Britain, populations peaked in the early 1970s, but have since declined by 68% overall, and about 90% in some regions. The RSPB lists the house sparrow's UK conservation status as red. In London, the house sparrow almost disappeared from the central city. The numbers of house sparrows in the Netherlands have dropped in half since the 1980s, so the house sparrow is even considered an endangered species. This status came to widespread attention after a female house sparrow, referred to as the "dominomus", was killed after knocking down dominoes arranged as part of an attempt to set a world record. These declines are not unprecedented, as similar reductions in population occurred when the internal combustion engine replaced horses in the 1920s and a major source of food in the form of grain spillage was lost.
Declines have been particularly apparent even in North America, where the house sparrow is invasive in some states. Introduced to Philadelphia initially in 1852 the house sparrow rapidly spread across the nation. However, the bird has largely disappeared from the city nowadays and overall, it is estimated to have declined in North America by 84% since 1966.
In South Asia, the house sparrow has largely vanished from major cities such as Karachi, Kolkata, Mumbai, New Delhi, and Lahore. In general, the house sparrow population has been on the decline in many Asian countries, and this is quite evident in India.
Various causes for the dramatic decreases in population have been proposed, including predation, in particular by Eurasian sparrowhawks, possibly facilitated by the elimination of bushes which the sparrows use to hide, electromagnetic radiation from mobile phones; and diseases such as avian malaria. A shortage of nesting sites caused by changes in urban building design is probably a factor, and conservation organisations have encouraged the use of special nest boxes for sparrows. A primary cause of the decline seems to be an insufficient supply of insect food for nestling sparrows. Declines in insect populations result from an increase of monoculture crops, the heavy use of pesticides, the replacement of native plants in cities with introduced plants and parking areas, and possibly the introduction of unleaded petrol, which produces toxic compounds such as methyl nitrite.
Protecting insect habitats on farms and planting native plants in cities benefit the house sparrow, as does establishing urban green spaces. To raise awareness of threats to the house sparrow, World Sparrow Day has been celebrated on 20 March across the world since 2010.
To promote the conservation of sparrow, in 2012, the house sparrow was declared as the state bird of Delhi.
Cultural associations
To many people across the world, the house sparrow is the most familiar wild animal and, because of its association with humans and familiarity, it is frequently used to represent the common and vulgar, or the lewd. One of the reasons for the introduction of house sparrows throughout the world was their association with the European homeland of many immigrants. Birds usually described later as sparrows are referred to in many works of ancient literature and religious texts in Europe and western Asia. These references may not always refer specifically to the house sparrow, or even to small, seed-eating birds, but later writers who were inspired by these texts often had the house sparrow in mind. In particular, sparrows were associated by the ancient Greeks with Aphrodite, the goddess of love, due to their perceived lustfulness, an association echoed by later writers such as Chaucer and Shakespeare. Jesus's use of "sparrows" as an example of divine providence in the Gospel of Matthew also inspired later references, such as that in Shakespeare's Hamlet and the Gospel hymn His Eye Is on the Sparrow.
G37
The house sparrow is very rarely represented in ancient Egyptian art, but an Egyptian hieroglyph is based on it. The sparrow hieroglyph had no phonetic value and was used as a determinative in words to indicate small, narrow, or bad. An alternative view is that the hieroglyph meant "a prolific man" or "the revolution of a year".
| Biology and health sciences | Passerida | null |
196206 | https://en.wikipedia.org/wiki/Longevity | Longevity | Longevity may refer to especially long-lived members of a population, whereas life expectancy is defined statistically as the average number of years remaining at a given age. For example, a population's life expectancy at birth is the same as the average age at death for all people born in the same year (in the case of cohorts).
Longevity studies may involve putative methods to extend life. Longevity has been a topic not only for the scientific community but also for writers of travel, science fiction, and utopian novels. The legendary fountain of youth appeared in the work of the Ancient Greek historian Herodotus.
There are difficulties in authenticating the longest human life span, owing to inaccurate or incomplete birth statistics. Fiction, legend, and folklore have proposed or claimed life spans in the past or future vastly longer than those verified by modern standards, and longevity narratives and unverified longevity claims frequently speak of their existence in the present.
A life annuity is a form of longevity insurance.
Life expectancy, as of 2010
Various factors contribute to an individual's longevity. Significant factors in life expectancy include gender, genetics, access to health care, hygiene, diet and nutrition, exercise, lifestyle, and crime rates. Below is a list of life expectancies in different types of countries:
Developed countries: 77–90 years (e.g. Canada: 81.29 years, 2010 est.)
Developing countries: 32–80 years (e.g. Mozambique: 41.37 years, 2010 est.)
Population longevities are increasing as life expectancies around the world grow:
Australia: 80 years in 2002, 81.72 years in 2010
France: 79.05 years in 2002, 81.09 years in 2010
Germany: 77.78 years in 2002, 79.41 years in 2010
Italy: 79.25 years in 2002, 80.33 years in 2010
Japan: 81.56 years in 2002, 82.84 years in 2010
Monaco: 79.12 years in 2002, 79.73 years in 2011
Spain: 79.06 years in 2002, 81.07 years in 2010
United Kingdom: 80 years in 2002, 81.73 years in 2010
United States: 77.4 years in 2002, 78.24 years in 2010
Long-lived individuals
The Gerontology Research Group validates current longevity records by modern standards, and maintains a list of supercentenarians; many other unvalidated longevity claims exist. Record-holding individuals include:
Eilif Philipsen (21 July 1682 – 20 June 1785, 102 years, 333 days): first person to reach the age of 100 (on 21 July 1782) and whose age could be validated.
Geert Adriaans Boomgaard (1788–1899, 110 years, 135 days): first person to reach the age of 110 (on September 21, 1898) and whose age could be validated.
Margaret Ann Neve, (18 May 1792 – 4 April 1903, 110 years, 346 days) the first validated female supercentenarian (on 18 May 1902).
Jeanne Calment (1875–1997, 122 years, 164 days): the oldest person in history whose age has been verified by modern documentation. This defines the modern human life span, which is set by the oldest documented individual who ever lived.
Sarah Knauss (1880–1999, 119 years, 97 days): the third oldest documented person in modern times and the oldest American.
Jiroemon Kimura (1897–2013, 116 years, 54 days): the oldest man in history whose age has been verified by modern documentation.
Kane Tanaka (1903–2022, 119 years, 107 days): the second oldest documented person in modern times and the oldest Japanese.
Major factors
Evidence-based studies indicate that longevity is based on two major factors: genetics and lifestyle.
Genetics
Twin studies have estimated that approximately 20-30% of the variation in human lifespan can be related to genetics, with the rest due to individual behaviors and environmental factors which can be modified. Although over 200 gene variants have been associated with longevity according to a US-Belgian-UK research database of human genetic variants these explain only a small fraction of the heritability.
Lymphoblastoid cell lines established from blood samples of centenarians have significantly higher activity of the DNA repair protein PARP (Poly ADP ribose polymerase) than cell lines from younger (20 to 70 year old) individuals. The lymphocytic cells of centenarians have characteristics typical of cells from young people, both in their capability of priming the mechanism of repair after sublethal oxidative DNA damage and in their PARP gene expression. These findings suggest that elevated PARP gene expression contributes to the longevity of centenarians, consistent with the DNA damage theory of aging.
In July 2020, scientists used public biological data on 1.75 m people with known lifespans overall and identified 10 genomic loci which appear to intrinsically influence healthspan, lifespan, and longevity – of which half have not been reported previously at genome-wide significance and most being associated with cardiovascular disease – and identified haem metabolism as a promising candidate for further research within the field. Their study suggests that high levels of iron in the blood likely reduce, and genes involved in metabolising iron likely increase healthy years of life in humans.
Lifestyle
Longevity is a highly plastic trait, and traits that influence its components respond to physical (static) environments and to wide-ranging life-style changes: physical exercise, dietary habits, living conditions, and pharmaceutical as well as nutritional interventions. A 2012 study found that even modest amounts of leisure time physical exercise can extend life expectancy by as much as 4.5 years.
Diet
As of 2021, there is no clinical evidence that any dietary practice contributes to human longevity. Although health can be influenced by diet, including the type of foods consumed, the amount of calories ingested, and the duration and frequency of fasting periods, there is no good clinical evidence that fasting promotes longevity in humans, .
Calorie restriction is a widely researched intervention to assess effects on aging, defined as a sustained reduction in dietary energy intake compared to the energy required for weight maintenance. To ensure metabolic homeostasis, the diet during calorie restriction must provide sufficient energy, micronutrients, and fiber. Some studies on rhesus monkeys showed that restricting calorie intake resulted in lifespan extension, while other animals studies did not detect a significant change. According to preliminary research in humans, there is little evidence that calorie restriction affects lifespan. There is a link between diet and obesity and consequent obesity-associated morbidity.
Biological pathways
Four well-studied biological pathways that are known to regulate aging, and whose modulation has been shown to influence longevity are Insulin/IGF-1, mechanistic target of rapamycin (mTOR), AMP-activating protein kinase (AMPK), and Sirtuin pathways.
Autophagy
Autophagy plays a pivotal role in healthspan and lifespan extension.
Change over time
In preindustrial times, deaths at young and middle age were more common than they are today. This is not due to genetics, but because of environmental factors such as disease, accidents, and malnutrition, especially since the former were not generally treatable with pre-20th-century medicine. Deaths from childbirth were common for women, and many children did not live past infancy. In addition, most people who did attain old age were likely to die quickly from the above-mentioned untreatable health problems. Despite this, there are several examples of pre-20th-century individuals attaining lifespans of 85 years or greater, including John Adams, Cato the Elder, Thomas Hobbes, Christopher Polhem, and Michelangelo. This was also true for poorer people like peasants or laborers. Genealogists will almost certainly find ancestors living to their 70s, 80s and even 90s several hundred years ago.
For example, an 1871 census in the UK (the first of its kind, but personal data from other censuses dates back to 1841 and numerical data back to 1801) found the average male life expectancy as being 44, but if infant mortality is subtracted, males who lived to adulthood averaged 75 years. The present life expectancy in the UK is 77 years for males and 81 for females, while the United States averages 74 for males and 80 for females.
Studies have shown that black American males have the shortest lifespans of any group of people in the US, averaging only 69 years (Asian-American females average the longest). This reflects overall poorer health and greater prevalence of heart disease, obesity, diabetes, and cancer among black American men.
Women normally outlive men. Theories for this include smaller bodies that place lesser strain on the heart (women have lower rates of cardiovascular disease) and a reduced tendency to engage in physically dangerous activities. Conversely, women are more likely to participate in health-promoting activities. The X chromosome also contains more genes related to the immune system, and women tend to mount a stronger immune response to pathogens than men. However, the idea that men have weaker immune systems due to the supposed immuno-suppressive actions of testosterone is unfounded.
There is debate as to whether the pursuit of longevity is a worthwhile health care goal. Bioethicist Ezekiel Emanuel, who is also one of the architects of ObamaCare, has argued that the pursuit of longevity via the compression of morbidity explanation is a "fantasy" and that longevity past age 75 should not be considered an end in itself. This has been challenged by neurosurgeon Miguel Faria, who states that life can be worthwhile in healthy old age, that the compression of morbidity is a real phenomenon, and that longevity should be pursued in association with quality of life. Faria has discussed how longevity in association with leading healthy lifestyles can lead to the postponement of senescence as well as happiness and wisdom in old age.
Naturally limited longevity
Most biological organisms have a naturally limited longevity due to aging, unlike a rare few that are considered biologically immortal.
Given that different species of animals and plants have different potentials for longevity, the disrepair accumulation theory of aging tries to explain how the potential for longevity of an organism is sometimes positively correlated to its structural complexity. It suggests that while biological complexity increases individual lifespan, it is counteracted in nature since the survivability of the overall species may be hindered when it results in a prolonged development process, which is an evolutionarily vulnerable state.
According to the antagonistic pleiotropy hypothesis, one of the reasons biological immortality is so rare is that certain categories of gene expression that are beneficial in youth become deleterious at an older age.
Myths and claims
Longevity myths are traditions about long-lived people (generally supercentenarians), either as individuals or groups of people, and practices that have been believed to confer longevity, but for which scientific evidence does not support the ages claimed or the reasons for the claims. A comparison and contrast of "longevity in antiquity" (such as the Sumerian King List, the genealogies of Genesis, and the Persian Shahnameh) with "longevity in historical times" (common-era cases through twentieth-century news reports) is elaborated in detail in Lucian Boia's 2004 book Forever Young: A Cultural History of Longevity from Antiquity to the Present and other sources.
After the death of Juan Ponce de León, Gonzalo Fernández de Oviedo y Valdés wrote in Historia General y Natural de las Indias (1535) that Ponce de León was looking for the waters of Bimini to cure his aging. Traditions that have been believed to confer greater human longevity also include alchemy, such as that attributed to Nicolas Flamel. In the modern era, the Okinawa diet has some reputation of linkage to exceptionally high ages.
Longevity claims may be subcategorized into four groups: "In late life, very old people often tend to advance their ages at the rate of about 17 years per decade .... Several celebrated super-centenarians (over 110 years) are believed to have been double lives (father and son, relations with the same names or successive bearers of a title) .... A number of instances have been commercially sponsored, while a fourth category of recent claims are those made for political ends ...." The estimate of 17 years per decade was corroborated by the 1901 and 1911 British censuses. Time magazine considered that, by the Soviet Union, longevity had been elevated to a state-supported "Methuselah cult".
Robert Ripley regularly reported supercentenarian claims in Ripley's Believe It or Not!, usually citing his own reputation as a fact-checker to claim reliability.
Non-human biological longevity
Longevity in other animals can shed light on the determinants of life expectancy in humans, especially when found in related mammals. However, important contributions to longevity research have been made by research in other species, ranging from yeast to flies to worms. In fact, some closely related species of vertebrates can have dramatically different life expectancies, demonstrating that relatively small genetic changes can have a dramatic impact on aging. For instance, Pacific Ocean rockfishes have widely varying lifespans. The species Sebastes minor lives a mere 11 years while its cousin Sebastes aleutianus can live for more than 2 centuries. Similarly, a chameleon, Furcifer labordi, is the current record holder for shortest lifespan among tetrapods, with only 4–5 months to live. By contrast, some of its relatives, such as Furcifer pardalis, have been found to live up to 6 years.
There are studies about aging-related characteristics of and aging in long-lived animals like various turtles and plants like Ginkgo biloba trees. They have identified potentially causal protective traits and suggest many of the species have "slow or [times of] negligible senescence" (or aging). The jellyfish T. dohrnii is biologically immortal and has been studied by comparative genomics.
Honey bees (Apis mellifera) are eusocial insects that display dramatic caste-specific differences in longevity. Queen bees live for an average of 1-2 years, compared to workers who live on average 15-38 days in summer and 150-200 days in winter. Worker honey bees with high amounts of flight experience exhibit increased DNA damage in flight muscle, as measured by elevated 8-Oxo-2'-deoxyguanosine, compared to bees with less flight experience. This increased DNA damage is likely due to an imbalance of pro- and anti-oxidants during flight-associated oxidative stress. Flight induced oxidative DNA damage appears to hasten senescence and reduce longevity in A. mellifera.
Examples of long-lived plants and animals
Currently living
Methuselah: over 4,850-year-old bristlecone pine in the White Mountains of California, the oldest currently living non-clonal tree.
Dead
WPN-114, "Prometheus": approximately 4,900 year-old (at time of tree-death) Pinus longaeva, located in Wheeler Peak, Nevada.
The quahog clam (Arctica islandica) is exceptionally long-lived, with a maximum recorded age of 507 years, the longest of any animal. Other clams of the species have been recorded as living up to 374 years.
Lamellibrachia luymesi, a deep-sea cold-seep tubeworm, is estimated to reach ages of over 250 years based on a model of its growth rates.
A bowhead whale killed in a hunt was found to be approximately 211 years old (possibly up to 245 years old), the longest-lived mammal known.
Possibly 250-million year-old bacteria, Bacillus permians, were revived from stasis after being found in sodium chloride crystals in a cavern in New Mexico.
Artificial animal longevity extension
Gene editing via CRISPR-Cas9 and other methods have significantly altered lifespans in animals.
| Biology and health sciences | Animal ontogeny | null |
196347 | https://en.wikipedia.org/wiki/Common%20pheasant | Common pheasant | The common pheasant (Phasianus colchicus), ring-necked pheasant, or blue-headed pheasant, a bird in the pheasant family (Phasianidae). The genus name comes from Latin phasianus 'pheasant'. The species name colchicus is Latin for 'of Colchis' (modern day Georgia), a country on the Black Sea where pheasants became known to Europeans. Although Phasianus was previously thought to be closely related to the genus Gallus, the genus of junglefowl and domesticated chickens, recent studies show that they are in different subfamilies, having diverged over 20 million years ago.
It is native to Asia, where it is widespread, and also the extreme southeast of Europe in the northern foothills of the Caucasus Mountains. It has been widely introduced elsewhere as a game bird. In parts of its range, mainly in places where none of its relatives occur such as in Europe, where it is naturalised, it is simply known as the "pheasant". Ring-necked pheasant is both the collective name for a number of subspecies and their intergrades that have white neck rings, and the name used for the species as a whole in North America.
It is a well-known gamebird, among those of more than regional importance perhaps the most widespread and ancient one in the whole world. The common pheasant is one of the world's most hunted birds; it has been introduced for that purpose to many regions, and is also common on game farms where it is commercially bred. The ring-necked subspecies group in particular are commonly bred and were introduced to many parts of the world; the game farm stock, though no distinct breeds have been developed yet, can be considered semi-domesticated. The ring-necked pheasant is the state bird of South Dakota, one of only two US state birds that is not a species native to the United States.
The green pheasant (P. versicolor) of Japan is sometimes considered a subspecies of the common pheasant. Though the species produce fertile hybrids wherever they coexist, this is simply a typical feature among fowl (Galloanseres), in which postzygotic isolating mechanisms are slight compared to most other birds. The species apparently have somewhat different ecological requirements and at least in its typical habitat, the green pheasant outcompetes the common pheasant. The introduction of the latter to Japan has therefore largely failed.
Description
There are many colour forms of the male common pheasant, ranging in colour from nearly white to almost black in some melanistic examples. These are due to captive breeding and hybridisation between subspecies and with the green pheasant, reinforced by continual releases of stock from varying sources to the wild. For example, the "ring-necked pheasants" common in Europe, North America and Australia do not pertain to any specific taxon, they rather represent a stereotyped hybrid swarm. Body weight can range from , with males averaging and females averaging . Wingspan ranges from .
The adult male common pheasant of the nominate subspecies Phasianus colchicus colchicus is in length with a long brown streaked black tail, accounting for almost of the total length. The body plumage is barred bright gold or fiery copper-red and chestnut-brown plumage with iridescent sheen of green and purple; but rump uniform is sometimes blue. The wing coverage is white or cream and black-barred markings are common on the tail. The head is bottle green with a small crest and distinctive red wattle. P. c. colchicus and some other races lack a white neck ring. Behind the face are two ear-tufts, that make the pheasant appear more alert.
The female (hen) and juveniles are much less showy, with a duller mottled brown plumage all over and measuring long including a tail of around . Juvenile birds have the appearance of the female with a shorter tail until young males begin to grow characteristic bright feathers on the breast, head and back at about 10 weeks after hatching.
The green pheasant (P. versicolor) is very similar, and hybridisation often makes the identity of individual farmed birds difficult to determine. Green pheasant males on average have a shorter tail than the common pheasant and have darker plumage that is uniformly bottle-green on the breast and belly; they always lack a neck ring. Green pheasant females are darker, with many black dots on the breast and belly.
In addition, various colour mutations are commonly encountered, mainly melanistic (black) and flavistic (isabelline or fawn) specimens. The former are rather commonly released in some areas and are named "tenebrosus pheasant" or simply "melanistic mutant".
Taxonomy and systematics
This species was first scientifically described by Carl Linnaeus in his landmark 1758 10th edition of Systema Naturae under its current scientific name. The common pheasant is distinct enough from any other species known to Linnaeus for a laconic [Phasianus] rufus, capîte caeruleo, "a red pheasant with blue head", to serve as entirely sufficient description. The bird had been extensively discussed before Linnaeus established binomial nomenclature so was already well-known. His sources are the Ornithologia of Ulisse Aldrovandi, Giovanni Pietro Olina's Uccelliera, John Ray's Synopsis methodica Avium & Piscium, and A Natural History of the Birds by Eleazar Albin. Therein—essentially the bulk of the ornithology textbooks of his day—the species is simply named "the pheasant" in the books' respective languages. Whereas in most other species, Linnaeus felt it warranted to cite plumage details from his sources, in the common pheasant's case he simply referred to the reason of the bird's fame: principum mensis dicatur. The type locality is given simply as "Africa, Asia".
However, the bird does not occur in Africa, except perhaps in Linnaeus's time in Mediterranean coastal areas where they might have been introduced during the Roman Empire. The type locality was later fixed to the Rioni River in western Georgia, known as Phasis to the Ancient Greeks. These birds, until the modern era, constituted the bulk of the introduced stock in parts of Europe that was not already present; the birds described by Linnaeus's sources, though typically belonging to such early introductions, would certainly have more alleles in common with the transcaucasian population than with others. The scientific name is Latin for "pheasant from Colchis", colchicus referring to the west of modern-day Georgia; the Ancient Greek term corresponding to the English "pheasant" is Phasianos ornis (Φασιανὸς ὂρνις), "bird of the river Phasis". Although Linnaeus included many Galliformes in his genus Phasianius such as the domestic chicken and its wild ancestor the red junglefowl, nowadays Gallus gallus, today only the common and the green pheasant are placed in this genus. As the latter was not known to Linnaeus in 1758, the common pheasant is treated as the type species of Phasianus.
In the US, common pheasants are widely known as "ring-necked pheasants". More colloquial North American names include "chinks" or, in Montana, "phezzens". In China, meanwhile, the species is properly called zhi ji (雉鸡), "pheasant-fowl", essentially implying the same as the English name "common pheasant". As elsewhere, P. colchicus is such a familiar bird in China that it is usually just referred to as shan ji (山雞), "mountain chicken", a Chinese term for pheasants in general.
As of 2005, it had the smallest known genome of all living amniotes, only 0.97 pg (970 million base pairs), roughly one-third of the human genome's size; however, the black-chinned hummingbird is the current holder of the smallest known amniote genome.
Subspecies
There are about 30 subspecies in five to eight groups. These can be identified by the male plumage, namely presence or absence of a white neck-ring and/or a white superciliary stripe, the colour of the uppertail (rump) and wing coverts, and the colour of crown, chest, upper back, and flank feathers. As noted above, introduced populations have mixed the alleles of various races by various amounts, differing according to the original stock used for introductions and what natural selection according to climate and habitat has made of that.
An investigation into the genetic relationships of subspecies suggested that the common pheasant originated from the forests of southeastern China. Initial divergence is thought to have occurred around 3.4 Mya. The lack of agreement between morphology-based subspecies delimitation and their genetic relationships is thought to be attributed to past isolation followed by more recent population mixing as the pheasant has expanded its range across the Palaearctic.
Sometimes this species is split into the Central Asian common and the East Asian ring-necked pheasants, roughly separated by the arid and high mountainous regions of Turkestan. However, while the western and eastern populations probably were entirely separate during the Zyryanka glaciation when deserts were more extensive, this separation was not long enough for actual speciation to occur. Today, the largest variety of colour patterns is found where the western and eastern populations mix, as is to be expected. Females usually cannot be identified even to subspecies group with certainty.
Many subspecies are in danger of disappearing due to hybridisation with introduced birds. The last black-necked pheasant (P. c. colchicus) population in Europe survives in Greece in the delta of the river Nestos, where in 2012 the population was estimated 100–250 individuals.
The subspecies groups, going from west to east, and some notable subspecies are:
Within a maximum clade credibility mDNA gene tree, the most basal group is the P. c. elegans-group of the Eastern Clade, diverging from the green pheasant during the Calabrian, and diversifying in Middle Pleistocene around 0.7 million years ago, with the groups of the Western Clade splitting off from those of the Eastern Clade about 0.59 million years ago. While the subspecies of the Western Clade are well geographically separated from each other, the subspecies of the Eastern Clade often show clinal variation and large areas of intergradation. For example, clines connect P. c. pallasi-karpowi-torquatus-takatsukasae within the P. c. torquatus group and P. c. kiangsuensis-alaschanicus-sohokhotensis-strauchi within the P. c. strauchi-vlangalii group, with the degree of expression of white collar and superciliary stripe in both cases decreasing from north to south. The isolated form P. c. hagenbecki is very close to P. c. pallasi in phenotype, and has been traditionally treated within the P. c. torquatus group until recently, when it was assigned in one study to the P. c. strauchi / P. c. vlangalii group. However, the origin of the corresponding feather samples as listed in GenBank is far away from the known distribution of subspecies P. c. hagenbecki, and the issue needs further clarification.
Ecology
Common pheasants are native to Asia and parts of Europe, their original range extending from the eastern Black Sea and the Caspian Sea to Manchuria, Siberia, Korea, Mainland China, and Taiwan. The birds are found in woodland, farmland, scrub, and wetlands. In its natural habitat the common pheasant lives in grassland near water with small copses of trees. Extensively cleared farmland is marginal habitat that cannot maintain self-sustaining populations for long.
Common pheasants are gregarious birds and outside the breeding season form loose flocks. However, captive bred common pheasants can show strong sexual segregation, in space and time, with sex differences in the use of feeding stations throughout the day. Wherever they are hunted they are always timid once they associate humans with danger, and will quickly retreat for safety after hearing the arrival of hunting parties in the area.
While common pheasants are able short-distance fliers, they prefer to run. If startled however, they can suddenly burst upwards at great speed, with a distinctive "whirring" wing sound and often giving kok kok kok calls to alert conspecifics. Their flight speed is only when cruising but when chased they can fly up to .
Nesting
Common pheasants nest solely on the ground in scrapes, lined with some grass and leaves, frequently under dense cover or a hedge. Occasionally they will nest in a haystack, or old nest left by other bird. They roost in sheltered trees at night. The males are polygynous as is typical for many Phasianidae, and are often accompanied by a harem of several females. Common pheasants produce a clutch of around 8–15 eggs, sometimes as many as 18, but usually 10 to 12; they are pale olive in colour, and laid over a 2–3 week period in April to June in the Northern Hemisphere. The incubation period is about 22–27 days. The chicks stay near the hen for several weeks, yet leave the nest when only a few hours old. After hatching they grow quickly, flying after 12–14 days, resembling adults by only 15 weeks of age.
They eat a wide variety of animal and vegetable type-food, like fruit, seeds, grain, mast, berries and leaves as well as a wide range of invertebrates, such as leatherjackets, ant eggs, wireworms, caterpillars, grasshoppers and other insects; with small vertebrates like lizards, field voles, small mammals and small birds occasionally taken.
European native
Southern Caucasian pheasants (P. c. colchicus) were common in Greece during the classical period and it is a widespread myth that the Greeks took pheasants to the Balkans when they colonised Colchis in the Caucasus. This colonization happened during the 6th century BC, but pheasant archaeological remains in the Balkans are much older dating to 6th millennium BC. This fact indicates that probably pheasants reached the area naturally. Additionally it seems that they had a continuous range in Turkey from the Sea of Marmara on the edge of the Balkans, across the northern shore of the country till Caucasus. The last remnants of the Balkan population survive in the Kotza-Orman riparian forest of Nestos, in Greece with an estimated population of 100–200 adult birds. In Bulgaria they were lost in the 1970s because they hybridised with introduced eastern subspecies.
Besides the Balkans the species lives in Europe in the area north of Caucasus where the local subspecies P.c.septentrionalis survives pure around the lower reaches of the Samur River. Reintroduction efforts in the rest of the north Caucasian range may include hybrid birds.
As an introduced species
Common pheasants can now be found across the globe due to their readiness to breed in captivity and the fact they can naturalise in many climates, but were known to be introduced in Europe, North America, Japan and New Zealand. Pheasants were hunted in their natural range by Stone Age humans just like the grouse, partridges, junglefowls and perhaps peafowls that inhabited Europe at that time. At least since the Roman Empire, the bird was extensively introduced in many places and has become a naturalized member at least of the European fauna. Introductions in the Southern Hemisphere have mostly failed, except where local Galliformes or their ecological equivalents are rare or absent.
The bird was naturalized in Great Britain around AD 1059, but may have been introduced by the Romano-British centuries earlier. It was the Caucasian subspecies mistakenly known as the 'Old English pheasant' rather than the Chinese ring-necked pheasants (torquatus) that were introduced to Britain. But it became extirpated from most of the isles in the early 17th century. There were further re-introductions of the 'white neck-ringed' variety in the 18th century. It was rediscovered as a game bird in the 1830s after being ignored for many years in an amalgam of forms. Since then it has been reared extensively by gamekeepers and was shot in season from 1 October to 31 January. Pheasants are well adapted to the British climate and breed naturally in the wild without human supervision in copses, heaths and commons.
By 1950 pheasants bred throughout the British Isles, although they were scarce in Ireland. Because around 30,000,000 pheasants are released each year on shooting estates, mainly in the Midlands and South of England, it is widespread in distribution, although most released birds survive less than a year in the wild. The Bohemian was most likely seen in North Norfolk. The Game & Wildlife Conservation Trust is researching the breeding success of reared pheasants and trying to find ways to improve this breeding success to reduce the demand to release as many reared pheasants and increase the wild population. As the original Caucasian stock all but disappeared during the Early Modern era, most 'dark-winged ringless' birds in the UK are actually descended from 'Chinese ring-necked' and 'green pheasant' hybrids, which were commonly used for rewilding.
North America
Common pheasants were introduced in North America in 1773, and have become well established throughout much of the Rocky Mountain states (Colorado, Idaho, Montana, Wyoming, etc.), the Midwest, the Plains states, as well as Canada and Mexico. In the southwest, they can even be seen south of the Rockies in Bosque del Apache National Wildlife Refuge south of Albuquerque, New Mexico. It is now most common on the Great Plains. Common pheasants have also been introduced to much of northwest Europe, the Hawaiian Islands, Chile, Uruguay, Peru, Argentina, Brazil, South Africa, New Zealand, and Australia including the island state of Tasmania and small offshore islands such as Rottnest Island off Western Australia.
Most common pheasants bagged in the United States are wild-born feral pheasants. In some states captive-reared and released birds make up much of the population.
Pheasant hunting is very popular in much of the US, especially in the Great Plains states, where a mix of farmland and native grasslands provides ideal habitat. South Dakota alone has an annual harvest of over 1 million birds a year by over 200,000 hunters.
Negative impacts on other birds
There are a number of negative effects of common pheasants on other game birds, including: nest parasitism, disease, aggression, and competition for resources. Nest parasitism, or brood parasitism, is common in pheasants because of their propensity to nest near other birds and the fact that nesting requirements are similar to those of other prairie birds and waterfowl that inhabit the same areas. This phenomenon has been observed in grey partridges; prairie chickens; several types of duck, rail, grouse, turkeys, and others. Effects of nest parasitism may include abandonment of nests with a high proportion of foreign eggs, lower hatching rates, and lower numbers of eggs laid by the host species. Pheasant eggs also have a shorter incubation time than many of their nestmates, which may result in the individual watching over the nest to abandon her own eggs after the pheasants hatch, thinking that the remaining eggs are not viable. Pheasants raised in other species' nests often imprint on their caretaker, which may result in them adopting atypical behaviour for their species. This is sometimes the cause of hybridisation of species as pheasants adopt the mating behaviour of their nest's host species.
Pheasants often compete with other native birds for resources. Studies have shown that they can lead to decreased populations of bobwhites and partridges due to habitat and food competition. Insects are a valuable food source for both pheasants and partridges and competition may lead to decreased populations of partridges. Pheasants may also introduce disease, such as blackhead, to native populations. While pheasants tolerate the infection well, other birds such as ruffed grouse, chukar, and grey partridge are highly susceptible. Pheasants also have a tendency to harass or kill other birds. One study noted that in pheasant vs. prairie chicken interactions, the pheasants were victorious 78% of the time.
Management strategies
A variety of management strategies have been suggested for areas that are home to species that are particularly threatened by pheasants, such as the prairie chickens and grey partridge. These strategies include mowing grass to decrease the nesting cover preferred by pheasants, decreasing pheasant roosting habitat, shooting pheasants in organised hunts, trapping and removing them from areas where there are high concentrations of birds of threatened species, and others.
Population change
While pheasant populations are not in any danger, they have been decreasing in the United States over the last 30 years, largely in agricultural areas. This is likely due to changes in farming practices, application of pesticides, habitat fragmentation, and increased predation due to changes in crops grown. Many crops beneficial for pheasants (such as barley) are not being farmed as much in favour of using the land for more lucrative crops, such as nut trees. Many of these new crops are detrimental to pheasant survival. Pheasants prefer to nest in areas of significant herbaceous cover, such as perennial grasses, so many agricultural areas are not conducive to nesting anymore. Pheasant hens also experience higher levels of predation in areas without patches of grassland.
In the United Kingdom, about 50 million pheasants reared in captivity are now released each summer, a number which has significantly increased since the 1980s. Most of these birds are shot during the open season (1 October to 1 February), and few survive for a year. The result is a wildly fluctuating population, from 50 million in July to less than 5 million in June.
As gamebirds
Common pheasants are bred to be hunted and are shot in great numbers in Europe, especially the UK, where they are shot on the traditional formal "driven shoot" principles, whereby paying guns have birds driven over them by beaters, and on smaller "rough shoots". The open season in the UK is 1 October – 1 February, under the Game Act 1831 (1 & 2 Will. 4. c. 32). Generally they are shot by hunters employing gun dogs to help find, flush and retrieve shot birds. Retrievers, spaniels and pointing breeds are used to hunt pheasants.
The doggerel "Up gets a guinea, bang goes a penny-halfpenny, and down comes a half a crown" reflects the expensive sport of 19th century driven shoots in Britain, when pheasants were often shot for sport, rather than as food. It was a popular royal pastime in Britain to shoot common pheasants. King George V shot over 1,000 pheasants out of a total bag of 3,937 over a six-day period in December 1913 during a competition with a friend; however, he did not do enough to beat him.
Common pheasants are traditionally a target of small game poachers in the UK. The Roald Dahl novel Danny the Champion of the World featured a poacher (and his son) who lived in the UK and illegally hunted common pheasants.
Pheasant farming is a common practice and is frequently done intensively, with serious adverse impacts on native species. Birds are supplied both to hunting preserves/estates and restaurants, with smaller numbers being available for home cooks.
The carcasses were often hung for a time to improve the meat by slight decomposition, as with most other game. Modern cookery generally uses moist roasting and farm-raised female birds. In the UK and US, game was making somewhat of a comeback in popular cooking and more pheasants than ever were being sold in supermarkets there in 2011. A major reason for this is consumer attitude shift from consumption of red meat to white meat.
| Biology and health sciences | Galliformes | null |
196493 | https://en.wikipedia.org/wiki/Non-coding%20RNA | Non-coding RNA | A non-coding RNA (ncRNA) is a functional RNA molecule that is not translated into a protein. The DNA sequence from which a functional non-coding RNA is transcribed is often called an RNA gene. Abundant and functionally important types of non-coding RNAs include transfer RNAs (tRNAs) and ribosomal RNAs (rRNAs), as well as small RNAs such as microRNAs, siRNAs, piRNAs, snoRNAs, snRNAs, exRNAs, scaRNAs and the long ncRNAs such as Xist and HOTAIR.
The number of non-coding RNAs within the human genome is unknown; however, recent transcriptomic and bioinformatic studies suggest that there are thousands of non-coding transcripts.
Many of the newly identified ncRNAs have unknown functions, if any.
There is no consensus on how much of non-coding transcription is functional: some believe most ncRNAs to be non-functional "junk RNA", spurious transcriptions, while others expect that many non-coding transcripts have functions to be discovered.
History and discovery
Nucleic acids were first discovered in 1868 by Friedrich Miescher, and by 1939, RNA had been implicated in protein synthesis. Two decades later, Francis Crick predicted a functional RNA component which mediated translation; he reasoned that RNA is better suited to base-pair with an mRNA transcript than a pure polypeptide.
The first non-coding RNA to be characterised was an alanine tRNA found in baker's yeast, its structure was published in 1965. To produce a purified alanine tRNA sample, Robert W. Holley et al. used 140kg of commercial baker's yeast to give just 1g of purified tRNAAla for analysis. The 80 nucleotide tRNA was sequenced by first being digested with Pancreatic ribonuclease (producing fragments ending in Cytosine or Uridine) and then with takadiastase ribonuclease Tl (producing fragments which finished with Guanosine). Chromatography and identification of the 5' and 3' ends then helped arrange the fragments to establish the RNA sequence. Of the three structures originally proposed for this tRNA, the 'cloverleaf' structure was independently proposed in several following publications. The cloverleaf secondary structure was finalised following X-ray crystallography analysis performed by two independent research groups in 1974.
Ribosomal RNA was next to be discovered, followed by URNA in the early 1980s. Since then, the discovery of new non-coding RNAs has continued with snoRNAs, Xist, CRISPR and many more. Recent notable additions include riboswitches and miRNA; the discovery of the RNAi mechanism associated with the latter earned Craig C. Mello and Andrew Fire the 2006 Nobel Prize in Physiology or Medicine.
Recent discoveries of ncRNAs have been achieved through both experimental and bioinformatic methods.
Biological roles
Noncoding RNAs belong to several groups and are involved in many cellular processes. These range from ncRNAs of central importance that are conserved across all or most cellular life through to more transient ncRNAs specific to one or a few closely related species. The more conserved ncRNAs are thought to be molecular fossils or relics from the last universal common ancestor and the RNA world, and their current roles remain mostly in regulation of information flow from DNA to protein.
In translation
Many of the conserved, essential and abundant ncRNAs are involved in translation. Ribonucleoprotein (RNP) particles called ribosomes are the 'factories' where translation takes place in the cell. The ribosome consists of more than 60% ribosomal RNA; these are made up of 3 ncRNAs in prokaryotes and 4 ncRNAs in eukaryotes. Ribosomal RNAs catalyse the translation of nucleotide sequences to protein. Another set of ncRNAs, Transfer RNAs, form an 'adaptor molecule' between mRNA and protein. The H/ACA box and C/D box snoRNAs are ncRNAs found in archaea and eukaryotes. RNase MRP is restricted to eukaryotes. Both groups of ncRNA are involved in the maturation of rRNA. The snoRNAs guide covalent modifications of rRNA, tRNA and snRNAs; RNase MRP cleaves the internal transcribed spacer 1 between 18S and 5.8S rRNAs. The ubiquitous ncRNA, RNase P, is an evolutionary relative of RNase MRP. RNase P matures tRNA sequences by generating mature 5'-ends of tRNAs through cleaving the 5'-leader elements of precursor-tRNAs. Another ubiquitous RNP called SRP recognizes and transports specific nascent proteins to the endoplasmic reticulum in eukaryotes and the plasma membrane in prokaryotes. In bacteria, Transfer-messenger RNA (tmRNA) is an RNP involved in rescuing stalled ribosomes, tagging incomplete polypeptides and promoting the degradation of aberrant mRNA.
In RNA splicing
In eukaryotes, the spliceosome performs the splicing reactions essential for removing intron sequences, this process is required for the formation of mature mRNA. The spliceosome is another RNP often known as the snRNP or tri-snRNP. There are two different forms of the spliceosome, the major and minor forms. The ncRNA components of the major spliceosome are U1, U2, U4, U5, and U6. The ncRNA components of the minor spliceosome are U11, U12, U5, U4atac and U6atac.
Another group of introns can catalyse their own removal from host transcripts; these are called self-splicing RNAs. There are two main groups of self-splicing RNAs: group I catalytic intron and group II catalytic intron. These ncRNAs catalyze their own excision from mRNA, tRNA and rRNA precursors in a wide range of organisms.
In mammals it has been found that snoRNAs can also regulate the alternative splicing of mRNA, for example snoRNA HBII-52 regulates the splicing of serotonin receptor 2C.
In nematodes, the SmY ncRNA appears to be involved in mRNA trans-splicing.
In DNA replication
Y RNAs are stem loops, necessary for DNA replication through interactions with chromatin and initiation proteins (including the origin recognition complex). They are also components of the Ro60 ribonucleoprotein particle which is a target of autoimmune antibodies in patients with systemic lupus erythematosus.
In gene regulation
The expression of many thousands of genes are regulated by ncRNAs. This regulation can occur in trans or in cis. There is increasing evidence that a special type of ncRNAs called enhancer RNAs, transcribed from the enhancer region of a gene, act to promote gene expression.
Trans-acting
In higher eukaryotes microRNAs regulate gene expression. A single miRNA can reduce the expression levels of hundreds of genes. The mechanism by which mature miRNA molecules act is through partial complementarity to one or more messenger RNA (mRNA) molecules, generally in 3' UTRs. The main function of miRNAs is to down-regulate gene expression.
The ncRNA RNase P has also been shown to influence gene expression. In the human nucleus, RNase P is required for the normal and efficient transcription of various ncRNAs transcribed by RNA polymerase III. These include tRNA, 5S rRNA, SRP RNA, and U6 snRNA genes. RNase P exerts its role in transcription through association with Pol III and chromatin of active tRNA and 5S rRNA genes.
It has been shown that 7SK RNA, a metazoan ncRNA, acts as a negative regulator of the RNA polymerase II elongation factor P-TEFb, and that this activity is influenced by stress response pathways.
The bacterial ncRNA, 6S RNA, specifically associates with RNA polymerase holoenzyme containing the sigma70 specificity factor. This interaction represses expression from a sigma70-dependent promoter during stationary phase.
Another bacterial ncRNA, OxyS RNA represses translation by binding to Shine-Dalgarno sequences thereby occluding ribosome binding. OxyS RNA is induced in response to oxidative stress in Escherichia coli.
The B2 RNA is a small noncoding RNA polymerase III transcript that represses mRNA transcription in response to heat shock in mouse
cells. B2 RNA inhibits transcription by binding to core Pol II. Through this interaction, B2 RNA assembles into preinitiation
complexes at the promoter and blocks RNA synthesis.
A recent study has shown that just the act of transcription of ncRNA sequence can have an influence on gene expression. RNA polymerase II transcription of ncRNAs is required for chromatin remodelling in the Schizosaccharomyces pombe. Chromatin is progressively converted to an open configuration, as several species of ncRNAs are transcribed.
Cis-acting
A number of ncRNAs are embedded in the 5' UTRs (Untranslated Regions) of protein coding genes and influence their expression in various ways. For example, a riboswitch can directly bind a small target molecule; the binding of the target affects the gene's activity.
RNA leader sequences are found upstream of the first gene of amino acid biosynthetic operons. These RNA elements form one of two possible structures in regions encoding very short peptide sequences that are rich in the end product amino acid of the operon. A terminator structure forms when there is an excess of the regulatory amino acid and ribosome movement over the leader transcript is not impeded. When there is a deficiency of the charged tRNA of the regulatory amino acid the ribosome translating the leader peptide stalls and the antiterminator structure forms. This allows RNA polymerase to transcribe the operon. Known RNA leaders are Histidine operon leader, Leucine operon leader, Threonine operon leader and the Tryptophan operon leader.
Iron response elements (IRE) are bound by iron response proteins (IRP). The IRE is found in UTRs of various mRNAs whose products are involved in iron metabolism. When iron concentration is low, IRPs bind the ferritin mRNA IRE leading to translation repression.
Internal ribosome entry sites (IRES) are RNA structures that allow for translation initiation in the middle of a mRNA sequence as part of the process of protein synthesis.
In genome defense
Piwi-interacting RNAs (piRNAs) expressed in
mammalian testes and somatic cells form RNA-protein complexes with Piwi proteins. These piRNA complexes (piRCs) have been linked to transcriptional gene silencing of retrotransposons and other genetic elements in germline cells, particularly those in spermatogenesis.
Clustered Regularly Interspaced Short Palindromic Repeats (CRISPR) are repeats found in the DNA of many bacteria and archaea. The repeats are separated by spacers of similar length. It has been demonstrated that these spacers can be derived from phage and subsequently help protect the cell from infection.
Chromosome structure
Telomerase is an RNP enzyme that adds specific DNA sequence repeats ("TTAGGG" in vertebrates) to telomeric regions, which are found at the ends of eukaryotic chromosomes. The telomeres contain condensed DNA material, giving stability to the chromosomes. The enzyme is a reverse transcriptase that carries Telomerase RNA, which is used as a template when it elongates telomeres, which are shortened after each replication cycle.
Xist (X-inactive-specific transcript) is a long ncRNA gene on the X chromosome of the placental mammals that acts as major effector of the X chromosome inactivation process forming Barr bodies. An antisense RNA, Tsix, is a negative regulator of Xist. X chromosomes lacking Tsix expression (and thus having high levels of Xist transcription) are inactivated more frequently than normal chromosomes. In drosophilids, which also use an XY sex-determination system, the roX (RNA on the X) RNAs are involved in dosage compensation. Both Xist and roX operate by epigenetic regulation of transcription through the recruitment of histone-modifying enzymes.
Bifunctional RNA
Bifunctional RNAs, or dual-function RNAs, are RNAs that have two distinct functions. The majority of the known bifunctional RNAs are mRNAs that encode both a protein and ncRNAs. However, a growing number of ncRNAs fall into two different ncRNA categories; e.g., H/ACA box snoRNA and miRNA.
Two well known examples of bifunctional RNAs are SgrS RNA and RNAIII. However, a handful of other bifunctional RNAs are known to exist (e.g., steroid receptor activator/SRA, VegT RNA,
Oskar RNA, ENOD40, p53 RNA SR1 RNA, and Spot 42 RNA.) Bifunctional RNAs were the subject of a 2011 special issue of Biochimie.
As a hormone
There is an important link between certain non-coding RNAs and the control of hormone-regulated pathways. In Drosophila, hormones such as ecdysone and juvenile hormone can promote the expression of certain miRNAs. Furthermore, this regulation occurs at distinct temporal points within Caenorhabditis elegans development. In mammals, miR-206 is a crucial regulator of estrogen-receptor-alpha.
Non-coding RNAs are crucial in the development of several endocrine organs, as well as in endocrine diseases such as diabetes mellitus. Specifically in the MCF-7 cell line, addition of 17β-estradiol increased global transcription of the noncoding RNAs called lncRNAs near estrogen-activated coding genes.
In pathogenic avoidance
C. elegans was shown to learn and inherit pathogenic avoidance after exposure to a single non-coding RNA of a bacterial pathogen.
Roles in disease
As with proteins, mutations or imbalances in the ncRNA repertoire within the body can cause a variety of diseases.
Cancer
Many ncRNAs show abnormal expression patterns in cancerous tissues. These include miRNAs, long mRNA-like ncRNAs, GAS5, SNORD50, telomerase RNA and Y RNAs. The miRNAs are involved in the large scale regulation of many protein coding genes, the Y RNAs are important for the initiation of DNA replication, telomerase RNA that serves as a primer for telomerase, an RNP that extends telomeric regions at chromosome ends (see telomeres and disease for more information). The direct function of the long mRNA-like ncRNAs is less clear.
Germline mutations in miR-16-1 and miR-15 primary precursors have been shown to be much more frequent in patients with chronic lymphocytic leukemia compared to control populations.
It has been suggested that a rare SNP (rs11614913) that overlaps hsa-mir-196a-2 has been found to be associated with non-small cell lung carcinoma. Likewise, a screen of 17 miRNAs that have been predicted to regulate a number of breast cancer associated genes found variations in the microRNAs miR-17 and miR-30c-1of patients; these patients were noncarriers of BRCA1 or BRCA2 mutations, lending the possibility that familial breast cancer may be caused by variation in these miRNAs.
The p53 tumor suppressor is arguably the most important agent in preventing tumor formation and progression. The p53 protein functions as a transcription factor with a crucial role in orchestrating the cellular stress response. In addition to its crucial role in cancer, p53 has been implicated in other diseases including diabetes, cell death after ischemia, and various neurodegenerative diseases such as Huntington, Parkinson, and Alzheimer. Studies have suggested that p53 expression is subject to regulation by non-coding RNA.
Another example of non-coding RNA dysregulated in cancer cells is the long non-coding RNA Linc00707. Linc00707 is upregulated and sponges miRNAs in human bone marrow-derived mesenchymal stem cells, gastric cancer or breast cancer, and thus promotes osteogenesis, contributes to hepatocellular carcinoma progression, promotes proliferation and metastasis, or indirectly regulates expression of proteins involved in cancer aggressiveness, respectively.
Prader–Willi syndrome
The deletion of the 48 copies of the C/D box snoRNA SNORD116 has been shown to be the primary cause of Prader–Willi syndrome. Prader–Willi is a developmental disorder associated with over-eating and learning difficulties. SNORD116 has potential target sites within a number of protein-coding genes, and could have a role in regulating alternative splicing.
Autism
The chromosomal locus containing the small nucleolar RNA SNORD115 gene cluster has been duplicated in approximately 5% of individuals with autistic traits. A mouse model engineered to have a duplication of the SNORD115 cluster displays autistic-like behaviour. A recent small study of post-mortem brain tissue demonstrated altered expression of long non-coding RNAs in the prefrontal cortex and cerebellum of autistic brains as compared to controls.
Cartilage–hair hypoplasia
Mutations within RNase MRP have been shown to cause cartilage–hair hypoplasia, a disease associated with an array of symptoms such as short stature, sparse hair, skeletal abnormalities and a suppressed immune system that is frequent among Amish and Finnish. The best characterised variant is an A-to-G transition at nucleotide 70 that is in a loop region two bases 5' of a conserved pseudoknot. However, many other mutations within RNase MRP also cause CHH.
Alzheimer's disease
The antisense RNA, BACE1-AS is transcribed from the opposite strand to BACE1 and is upregulated in patients with Alzheimer's disease. BACE1-AS regulates the expression of BACE1 by increasing BACE1 mRNA stability and generating additional BACE1 through a post-transcriptional feed-forward mechanism. By the same mechanism it also raises concentrations of beta amyloid, the main constituent of senile plaques. BACE1-AS concentrations are elevated in subjects with Alzheimer's disease and in amyloid precursor protein transgenic mice.
miR-96 and hearing loss
Variation within the seed region of mature miR-96 has been associated with autosomal dominant, progressive hearing loss in humans and mice. The homozygous mutant mice were profoundly deaf, showing no cochlear responses. Heterozygous mice and humans progressively lose the ability to hear.
Mitochondrial transfer RNAs
A number of mutations within mitochondrial tRNAs have been linked to diseases such as MELAS syndrome, MERRF syndrome, and chronic progressive external ophthalmoplegia.
Distinction between functional RNA (fRNA) and ncRNA
Scientists have started to distinguish functional RNA (fRNA) from ncRNA, to describe regions functional at the RNA level that may or may not be stand-alone RNA transcripts. This implies that fRNA (such as riboswitches, SECIS elements, and other cis-regulatory regions) is not ncRNA. Yet fRNA could also include mRNA, as this is RNA coding for protein, and hence is functional. Additionally artificially evolved RNAs also fall under the fRNA umbrella term. Some publications state that ncRNA and fRNA are nearly synonymous, however others have pointed out that a large proportion of annotated ncRNAs likely have no function. It also has been suggested to simply use the term RNA, since the distinction from a protein coding RNA (messenger RNA) is already given by the qualifier mRNA. This eliminates the ambiguity when addressing a gene "encoding a non-coding" RNA. Besides, there may be a number of ncRNAs that are misannoted in published literature and datasets.
| Biology and health sciences | Molecular biology | Biology |
196523 | https://en.wikipedia.org/wiki/Toxoplasma%20gondii | Toxoplasma gondii | Toxoplasma gondii () is a species of parasitic alveolate that causes toxoplasmosis. Found worldwide, T. gondii is capable of infecting virtually all warm-blooded animals, but felids are the only known definitive hosts in which the parasite may undergo sexual reproduction.
In rodents, T. gondii alters behavior in ways that increase the rodents' chances of being preyed upon by felids. Support for this "manipulation hypothesis" stems from studies showing that T. gondii-infected rats have a decreased aversion to cat urine while infection in mice lowers general anxiety, increases explorative behaviors and increases a loss of aversion to predators in general. Because cats are one of the only hosts within which T. gondii can sexually reproduce, such behavioral manipulations are thought to be evolutionary adaptations that increase the parasite's reproductive success since rodents that do not avoid cat habitations will more likely become cat prey. The primary mechanisms of T. gondii–induced behavioral changes in rodents occur through epigenetic remodeling in neurons that govern the relevant behaviors (e.g. hypomethylation of arginine vasopressin-related genes in the medial amygdala, which greatly decrease predator aversion).
In humans, particularly infants and those with weakened immunity, T. gondii infection is generally asymptomatic but may lead to a serious case of toxoplasmosis. T. gondii can initially cause mild, flu-like symptoms in the first few weeks following exposure, but otherwise, healthy human adults are asymptomatic. This asymptomatic state of infection is referred to as a latent infection, and it has been associated with numerous subtle behavioral, psychiatric, and personality alterations in humans. Behavioral changes observed between infected and non-infected humans include a decreased aversion to cat urine (but with divergent trajectories by gender) and an increased risk of schizophrenia. Preliminary evidence has suggested that T. gondii infection may induce some of the same alterations in the human brain as those observed in rodents. Many of these associations have been strongly debated and newer studies have found them to be weak, concluding:
However, there is evidence that T. Gondii may cause suicidal ideation and suicide in humans.
T. gondii is one of the most common parasites in developed countries; serological studies estimate that up to 50% of the global population has been exposed to, and may be chronically infected with, T. gondii; although infection rates differ significantly from country to country. Estimates have shown the highest IgG seroprevalence to be in Ethiopia, at 64.2%, as of 2018.
Structure
T. gondii contains organelles called rhoptries and micronemes. They contain proteins for invasion and effectors for manipulating the hosts immune response. To inject them into host cells, T. gondii uses the apical complex located in the tip of the cell to puncture the host membrane and discharge the contents of these organelles. The two microtubulins and their associated proteins inside the conoid of the apical complex facilitate this by organizing and docking the rhoptries to the complex. The mechanism underlying this is yet to be discovered fully but the roles of the microtubulins and the four associated proteins have been identified.
Life cycle
The life cycle of T. gondii may be broadly summarized into two components: a sexual component that occurs only within cats (felids, wild or domestic), and an asexual component that can occur within virtually all warm-blooded animals, including humans, cats, and birds. Because T. gondii can sexually reproduce only within cats, cats are therefore the definitive host of T. gondii. All other hosts – in which only asexual reproduction can occur – are intermediate hosts.
Sexual reproduction in the feline definitive host
When a feline is infected with T. gondii (e.g. by consuming an infected mouse carrying the parasite's tissue cysts), the parasite survives passage through the stomach, eventually infecting epithelial cells of the cat's small intestine. Inside these intestinal cells, the parasites undergo sexual development and reproduction, producing millions of thick-walled, zygote-containing cysts known as oocysts. Felines are the only definitive host because they lack expression of the enzyme delta-6-desaturase (D6D) in their intestine. This enzyme converts linoleic acid; the absence of expression allows systemic linoleic acid accumulation. Recent findings showed that this excess of linoleic acid is essential for T. gondii sexual reproduction.
Feline shedding of oocysts
Infected epithelial cells eventually rupture and release oocysts into the intestinal lumen, whereupon they are shed in the cat's feces. Oocysts can then spread to soil, water, food, or anything potentially contaminated with the feces. Highly resilient, oocysts can survive and remain infective for many months in cold and dry climates.
Ingestion of oocysts by humans or other warm-blooded animals is one of the common routes of infection. Humans can be exposed to oocysts by, for example, consuming unwashed vegetables or contaminated water, or by handling the feces (litter) of an infected cat. Although cats can also be infected by ingesting oocysts, they are much less sensitive to oocyst infection than are intermediate hosts.
Initial infection of the intermediate host
Intermediate hosts found include pigs, chickens, goats, sheep and Macropus rufus by Moré et al. 2010. Cattle and horses are resistant and thought to be incapable of significant infection. T. gondii is considered to have three stages of infection; the tachyzoite stage of rapid division, the bradyzoite stage of slow division within tissue cysts, and the oocyst environmental stage. Tachyzoites are also known as "tachyzoic merozoites" and bradyzoites as "bradyzoic merozoites". When an oocyst or tissue cyst is ingested by a human or other warm-blooded animal, the resilient cyst wall is dissolved by proteolytic enzymes in the stomach and small intestine, freeing sporozoites from within the oocyst. The parasites first invade cells in and surrounding the intestinal epithelium, and inside these cells, the parasites differentiate into tachyzoites, the motile and quickly multiplying cellular stage of T. gondii. Tissue cysts in tissues such as brain and muscle tissue, form about 7–10 days after initial infection. Although severe infection of M. rufus has been observed it is unknown whether this is common.
Asexual reproduction in the intermediate host
Inside host cells, the tachyzoites replicate inside specialized vacuoles (called the parasitophorous vacuoles) created from host cell membrane during invasion into the cell. Tachyzoites multiply inside this vacuole until the host cell dies and ruptures, releasing and spreading the tachyzoites via the bloodstream to all organs and tissues of the body, including the brain.
Growth in tissue culture
The parasite can be easily grown in monolayers of mammalian cells maintained in vitro in tissue culture. It readily invades and multiplies in a wide variety of fibroblast and monocyte cell lines. In infected cultures, the parasite rapidly multiplies and thousands of tachyzoites break out of infected cells and enter adjacent cells, destroying the monolayer in due course. New monolayers can then be infected using a drop of this infected culture fluid and the parasite indefinitely maintained without the need of animals.
Formation of tissue cysts
Following the initial period of infection characterized by tachyzoite proliferation throughout the body, pressure from the host's immune system causes T. gondii tachyzoites to convert into bradyzoites, the semidormant, slowly dividing cellular stage of the parasite. Inside host cells, clusters of these bradyzoites are known as tissue cysts. The cyst wall is formed by the parasitophorous vacuole membrane. Although bradyzoite-containing tissue cysts can form in virtually any organ, tissue cysts predominantly form and persist in the brain, the eyes, and striated muscle (including the heart). However, specific tissue tropisms can vary between intermediate host species; in pigs, the majority of tissue cysts are found in muscle tissue, whereas in mice, the majority of cysts are found in the brain.
Cysts usually range in size between five and 50 μm in diameter, (with 50 μm being about two-thirds the width of the average human hair).
Consumption of tissue cysts in meat is one of the primary means of T. gondii infection, both for humans and for meat-eating, warm-blooded animals. Humans consume tissue cysts when eating raw or undercooked meat (particularly pork and lamb). Tissue cyst consumption is also the primary means by which cats are infected.
An exhibit at the San Diego Natural History Museum states urban runoff with cat feces transports Toxoplasma gondii into the ocean, which can kill sea otters.
Chronic infection
Tissue cysts can be maintained in host tissue for the lifetime of the animal. However, the perpetual presence of cysts appears to be due to a periodic process of cyst rupturing and re-encysting, rather than a perpetual lifespan of individual cysts or bradyzoites. At any given time in a chronically infected host, a very small percentage of cysts are rupturing, although the exact cause of this tissue cyst rupture is, as of 2010, not yet known.
Theoretically, T. gondii can be passed between intermediate hosts indefinitely via a cycle of consumption of tissue cysts in meat. However, the parasite's life cycle begins and completes only when the parasite is passed to a feline host, the only host within which the parasite can again undergo sexual development and reproduction.
Population structure in the wild
In 2006, researchers reviewed evidence that T. gondii has an unusual population structure dominated by three clonal lineages called Types I, II and III that occur in North America and Europe, despite the occurrence of a sexual phase in its life cycle. They estimated that a common ancestor existed about 10,000 years ago. Authors of a subsequent and larger study on 196 isolates from diverse sources including T. gondii in the bald eagle, gray wolf, Arctic fox and sea otter, also found that T. gondii strains infecting North American wildlife have limited genetic diversity with the occurrence of only a few major clonal types. They found that 85% of strains in North America were of one of three widespread genotypes II, III and Type 12. Thus T. gondii has retained the capability for sex in North America over many generations, producing largely clonal populations, and matings have generated little genetic diversity.
Cellular stages
During different periods of its life cycle, individual parasites convert into various cellular stages, with each stage characterized by a distinct cellular morphology, biochemistry, and behavior. These stages include the tachyzoites, merozoites, bradyzoites (found in tissue cysts), and sporozoites (found in oocysts).
Some stages are motile and some calcium-dependent protein kinases (s) are involved in this parasite's motility. Gaji et al. 2015 find is required to begin the action of motility because it phosphorylates T. gondiis myosin A (). TgCDPK3 is the functional orthologue of CDPK1 in this parasite.
Tachyzoites
Motile, and quickly multiplying, tachyzoites are responsible for expanding the population of the parasite in the host. When a host consumes a tissue cyst (containing bradyzoites) or an oocyst (containing sporozoites), the bradyzoites or sporozoites stage-convert into tachyzoites upon infecting the intestinal epithelium of the host. During the initial acute period of infection, tachyzoites spread throughout the body via the blood stream. During the later, latent (chronic) stages of infection, tachyzoites stage-convert to bradyzoites to form tissue cysts. To survive in the host, tachyzoites manipulate the immune response by injecting the contents of rhoptries into host cells. This seems to be vital for their survival, as knock-out strains of T. gondii are unable to inject hosts with rhoptries have been shown to be avirulent in vivo.
Merozoites
Like tachyzoites, merozoites divide quickly and are responsible for expanding the population of the parasite inside the cat's intestine before sexual reproduction. When a feline definitive host consumes a tissue cyst (containing bradyzoites), bradyzoites convert into merozoites inside intestinal epithelial cells. Following a brief period of rapid population growth in the intestinal epithelium, merozoites convert into the noninfectious sexual stages of the parasite to undergo sexual reproduction, eventually resulting in zygote-containing oocysts.
Studying the sexual phases of the T. gondii life cycle remains challenging and determining the precise triggers and molecular mechanisms governing this developmental program remains an ongoing area of research. Major challenges associated with the ability to cultivate presexual and sexual stages of T. gondii in vitro have limited our understanding of this developmental program and how it is triggered by the parasite in response to the infection of the cat. Multiple studies , revealed distinct differences in the transcriptomes of the asexual and sexual stages of T. gondii. Additionally, metabolic disparities within the feline host have been identified as key factors influencing the transition to sexual stages. However, linking gene expression patterns to stage transitions and deciphering the genetic triggers driving the switch from asexual to sexual development remain unresolved.
Important recent advancements in the field have shed new light on the regulatory mechanisms governing sexual development in T. gondii. Farhat and colleagues showed that chromatin modifiers MORC and HDAC3 play critical roles in silencing sexual development-specific genes. In MORC-depleted parasites, a broad activation of sexual gene expression was observed. In a later study, it was suggested that MORC-depleted parasites have disrupted sub-telomeric gene silencing. The disorganization in telomeres may have led to the misregulation of sexual development.
Moreover, the discovery of specific transcription factors essential for sexual commitment has provided invaluable insights into the intricate regulatory network orchestrating stage specificity in T. gondii. Multiple parasite transcription factors have been identified as critical suppressors of presexual development, permitting the study of presexual stages and opening new avenues for using genetics to drive the full sexual cycle in vitro. Specifically, the depletion of AP2XI-2 and AP2XII-1 in T. gondii induces merozoite-specific gene expression, raising the possibility for cultivating T. gondii sexual development in laboratory settings.
Crucial questions still persist regarding the genetic determinants that dictate whether parasites develop into macrogametes or microgametes. The development of new molecular and genomic approaches, such as single-cell transcriptomics and proteomics, should be useful to those in the field working towards unraveling the molecular intricacies of this process.
Bradyzoites
Bradyzoites are the slowly dividing stage of the parasite that make up tissue cysts. When an uninfected host consumes a tissue cyst, bradyzoites released from the cyst infect intestinal epithelial cells before converting to the proliferative tachyzoite stage. Following the initial period of proliferation throughout the host body, tachyzoites then convert back to bradyzoites, which reproduce inside host cells to form tissue cysts in the new host.
Sporozoites
Sporozoites are the stage of the parasite residing within oocysts. When a human or other warm-blooded host consumes an oocyst, sporozoites are released from it, infecting epithelial cells before converting to the proliferative tachyzoite stage.
Immune response
Initially, a T. gondii infection stimulates production of IL-2 and IFN-γ by the innate immune system. Continuous IFN-γ production is necessary for control of both acute and chronic T. gondii infection. These two cytokines elicit a CD4+ and CD8+ T-cell mediated immune response. Thus, T-cells play a central role in immunity against Toxoplasma infection. T-cells recognize Toxoplasma antigens that are presented to them by the body's own Major Histocompatibility Complex (MHC) molecules. The specific genetic sequence of a given MHC molecule differs dramatically between individuals, which is why these molecules are involved in transplant rejection. Individuals carrying certain genetic sequences of MHC molecules are much more likely to be infected with Toxoplasma. One study of >1600 individuals found that Toxoplasma infection was especially common among people who expressed certain MHC alleles (HLA-B*08:01, HLA-C*04:01, HLA-DRB 03:01, HLA-DQA*05:01 and HLA-DQB*02:01).
IL-12 is produced during T. gondii infection to activate natural killer (NK) cells. Tryptophan is an essential amino acid for T. gondii, which it scavenges from host cells. IFN-γ induces the activation of indole-amine-2,3-dioxygenase (IDO) and tryptophan-2,3-dioxygenase (TDO), two enzymes that are responsible for the degradation of tryptophan. Immune pressure eventually leads the parasite to form cysts that normally are deposited in the muscles and in the brain of the hosts.
Immune response and behavior alterations
The IFN-γ-mediated activation of IDO and TDO is an evolutionary mechanism that serves to starve the parasite, but it can result in depletion of tryptophan in the brain of the host. IDO and TDO degrade tryptophan to N-formylkynurenine. Administration of L-kynurenine is capable of inducing depressive-like behavior in mice. T. gondii infection has been demonstrated to increase the levels of kynurenic acid (KYNA) in the brains of infected mice and in the brain of schizophrenic persons. Low levels of tryptophan and serotonin in the brain were already associated with depression.
Risk factors for human infection
The following have been identified as being risk factors for T. gondii infection in humans and warm-blooded animals:
by consuming raw or undercooked meat containing T. gondii tissue cysts. The most common threat to citizens in the United States is from eating raw or undercooked pork.
by ingesting water, soil, vegetables, or anything contaminated with oocysts shed in the feces of an infected animal. Cat fecal matter is particularly dangerous: Just one cyst consumed by a cat can result in thousands of oocysts. This is why physicians recommend pregnant or ill persons do not clean the cat's litter box at home. These oocysts are resilient to harsh environmental conditions and can survive over a year in contaminated soil.
from a blood transfusion or organ transplant
from transplacental transmission from mother to fetus, particularly when T. gondii is contracted during pregnancy
from drinking unpasteurized goat milk
from raw and treated sewage and bivalve shellfish contaminated by treated sewage
A common argument in the debate about whether cat ownership is ethical involves the question of T. gondii transmission to humans. Even though "living in a household with a cat that used a litter box was strongly associated with infection," and that living with several kittens or any cat under one year of age has some significance, several other studies claim to have shown that living in a household with a cat is not a significant risk factor for T. gondii infection.
Specific vectors for transmission may also differ based on geographic location. "The seawater in California is thought to be contaminated by T. gondii oocysts that originate from cat feces, survive or bypass sewage treatment, and travel to the coast through river systems. T. gondii has been identified in a California mussel by polymerase chain reaction and DNA sequencing. In light of the potential presence of T. gondii, pregnant women and immunosuppressed persons should be aware of this potential risk associated with eating raw oysters, mussels, and clams."
In warm-blooded animals, such as brown rats, sheep, and dogs, T. gondii has also been shown to be sexually transmitted. Although T. gondii can infect, be transmitted by, and asexually reproduce within humans and virtually all other warm-blooded animals, the parasite can sexually reproduce only within the intestines of members of the cat family (felids). Felids are therefore the definitive hosts of T. gondii; all other hosts (such as human or other mammals) are intermediate hosts.
Preventing infection
The following precautions are recommended to prevent or greatly reduce the chances of becoming infected with T. gondii. This information has been adapted from the websites of United States Centers for Disease Control and Prevention and the Mayo Clinic.
From food
Basic food-handling safety practices can prevent or reduce the chances of becoming infected with T. gondii, such as washing unwashed fruits and vegetables, and avoiding raw or undercooked meat, poultry, and seafood. Other unsafe practices such as drinking unpasteurized milk or untreated water can increase odds of infection. As T. gondii is commonly transmitted through ingesting microscopic cysts in the tissues of infected animals, meat that is not prepared to destroy these presents a risk of infection. Freezing meat for several days at subzero temperatures (0 °F or −18 °C) before cooking may break down all cysts, as they rarely survive these temperatures. During cooking, whole cuts of red meat should be cooked to an internal temperature of at least 145 °F (63 °C). Medium rare meat is generally cooked between 130 and 140 °F (55 and 60 °C), so cooking meat to at least medium is recommended. After cooking, a rest period of 3 min should be allowed before consumption. However, ground meat should be cooked to an internal temperature of at least 160 °F (71 °C) with no rest period. All poultry should be cooked to an internal temperature of at least 165 °F (74 °C). After cooking, a rest period of 3 min should be allowed before consumption.
From environment
Oocysts in cat feces take at least a day to sporulate (to become infectious after they are shed), so disposing of cat litter daily greatly reduces the chance of infectious oocysts developing. As these can spread and survive in the environment for months, humans should wear gloves when gardening or working with soil, and should wash their hands promptly after disposing of cat litter. These precautions apply to outdoor sandboxes/play sand pits, which should be covered when not in use. Cat feces should never be flushed down a toilet.
Pregnant women are at higher risk of transmitting the parasite to their unborn child and immunocompromised people of acquiring a lingering infection. Because of this, they should not change or handle cat litter boxes. Ideally, cats should be kept indoors and fed only food that has low to no risk of carrying oocysts, such as commercial cat food or well-cooked table food.
Vaccination
No approved human vaccine exists against Toxoplasma gondii. Research on human vaccines is ongoing.
For sheep, an approved live vaccine sold as Toxovax (from MSD Animal Health) provides lifetime protection.
There is currently no commercially available vaccine to prevent T. gondii infection in cats. However, research into feline vaccines for toxoplasmosis is ongoing, with several candidates showing positive results in clinical trials.
Treatment
In humans, active toxoplasmosis can be treated with a combination of drugs such as pyrimethamine and sulfadiazine, plus folinic acid. Immune-compromised patients may need continuous treatment until/unless their immune system is restored.
Environmental effects
In many parts of the world, where there are high populations of feral cats, there is an increased risk to the native wildlife due to increased infection of Toxoplasma gondii. It has been found that the serum concentrations of T. gondii in the wildlife population were increased where there are high amounts of cat populations. This creates a dangerous environment for organisms that have not evolved in cohabitation with felines and their contributing parasites.
Impact on marine species
Cetaceans
Toxoplasmosis has been implicated in the deaths of various cetaceans species, such as the critically endangered Māui dolphin and Hector's dolphin found in New Zealand. With only 54 Māui dolphins over the age of one remaining, T. gondii is considered a significant human-caused threat to the dolphins’ populations. Fatal cases of T. gondii have also been confirmed among spinner dolphins off the coast of Hawaii, among bottlenose dolphins, Risso's dolphins, and striped dolphins along the Mediterranean coast, among Indo-Pacific humpback dolphins in Australia, and again among free-ranging bottlenose dolphins in Brazil.
A 2011 study of 161 Pacific Northwest marine mammals ranging from a sperm whale to harbor porpoises that had either become stranded or died found that 42 percent tested positive for both T. gondii and S. neurona. Approximately 14 per cent of the western Arctic beluga whale population is believed to asymptomatically carry T. gondii with a few deaths attributed to the infection.
Minks and otters
Toxoplasmosis is one of the contributing factors toward mortality in southern sea otters, especially in areas where there is large urban run-off. In their natural habitats, sea otters control sea urchin populations and, thus indirectly, control sea kelp forests. By enabling the growth of sea kelp, other marine populations are protected as well as CO2 emissions are reduced due to the kelp's ability to absorb atmospheric carbon. An examination on 105 beachcast otters revealed that 38.1% had parasitic infections, and 28% of said infections had resulted in protozoal meningoencephalitis deaths. Toxoplasma gondii was found to be the root cause in 16.2% of these deaths, while 6.7% of the deaths were due to a closely related protozoan parasite known as Sarcocystis neurona.
Minks, being semiaquatic, are also susceptible to infection and being antibody-positive toward T. gondii. Minks can follow a similar diet as otters and feasts on crustaceans, fish, and invertebrates, thus the transmission route follows a similar pattern to otters. Because of the mink's ability to transverse land more frequently, and often seen as an invasive species itself, minks are a bigger threat in transporting T. gondii to other mammalian species, rather than otters who have a more restrictive breadth.
Other marine mammals
T. gondii has killed at least twelve endangered Hawaiian monk seals. There has been a documented fatal case in a West Indian manatee.
Black-footed penguins
Although under-studied, penguin populations, especially those that share an environment with the human population, are at-risk due to parasite infections, mainly Toxoplasmosis gondii. The main subspecies of penguins found to be infected by T. gondii include wild Magellanic and Galapagos penguins, as well as blue and African penguins in captivity. In one study, 57 (43.2%) of 132 serum samples of Magellanic penguins were found to have T. gondii. The island that the penguin is located, Magdalena Island, is known to have no cat populations, but a very frequent human population, indicating the possibility of transmission.
Histopathology
Examination of black-footed penguins with toxoplasmosis reveals hepatomegaly, splenomegaly, cranial hemorrhage, and necrotic kidneys. Alveolar and hepatic tissue presents a high number of immune cells such as macrophages containing tachyzoites of T. gondii. Histopathological features in other animals affected with toxoplasmosis had tachyzoites in eye structures such as the retina which lead to blindness.
Water transmission
The transmission of oocysts has been unknown, even though there have many documented cases of infection in marine species. Researchers have found that the oocytes of T. gondii can survive in seawater for at least six months, with the amount of salt concentration not affecting its life cycle. There have been no studies on the ability of T. gondii oocysts life cycle within freshwater environments, although infections are still present. One possible hypothesis of transmission is via amoeba species, particularly Acanthamoeba spp., a species that is found in all water environments (fresh, brackish, and full-strength seawater). Normally, amoebas function as a natural filter, phagocytizing nutrients and bacteria found within the water. Some pathogens have used this to their advantage, however, and evolved to be able to avoid being broken down and, thus, survive encased in the amoeba – this includes Holosporaceae, Pseudomonaceae, Burkholderiacceae, among others. Overall, this aids the pathogen in transportation but, also, protection from drugs and sterilizers that would, otherwise, cause death in the pathogen. Studies have shown that T. gondii oocysts can live within amoebas after being engulfed for at least 14 days without significant obliteration of the parasite. The ability of the microorganism to survive in vitro is dependent on the microorganism itself, but there are a few overarching mechanisms present. T. gondii oocysts have been found to resist an acidic pH and, thus, are protected by the acidification found in endocytic vacuoles and lysosomes. Phagocytosis further increases with the carbohydrate-rich surface membrane located on the amoebae. The pathogen can be released either by lysis of the amoebae or by exocytosis, but this is understudied
Impact on wild birds
Almost all species of birds that have been tested for Toxoplasma gondii have shown to be positive. The only bird species not reported with clinical symptoms of toxoplasmosis would be wild ducks, and there has only been one report found on domesticated ducks occurring in 1962. Species with resistance toward T. gondii include domestic turkeys, owls, red tail hawks, and sparrows, depending on the strain of T. gondii. T. gondii is considerably more severe in pigeons, particularly crown pigeons, ornamental pigeons, and pigeons originating from Australia and New Zealand. Typical onset is quick and usually results in death. Those that do survive often have chronic conditions of encephalitis and neuritis. Similarly, canaries are observed to be just as severe as pigeons, but the clinical symptoms are more abnormal when compared to other species. Most of the infection affects the eye, causing blindness, choroidal lesions, conjunctivitis, atrophy of the eye, blepharitis, and chorioretinitis Most of the time, the infection leads to death. Research by Michael Grigg, chief of the molecular parasitology unit at the National Institute of Allergy and Infectious Diseases, found that more than one half of dead raptors and more than one third of dead seabirds examined had the T. gondii parasite.
Current environmental efforts
Urbanization and global warming are extremely influential in the transmission of T. gondii. Temperature and humidity are huge factors in the sporulation stage: low humidity is always fatal to the oocysts, and they are also vulnerable to extreme temperatures. Rainfall is also an important factor for survival of waterborne pathogens. Because increased rainfall directly increases the flow rate in rivers, the amount of flow into coastal areas is increased as well. This can spread waterborne pathogens over wide areas.
There is no effective vaccine for T. gondii, and research on a live vaccine is ongoing. Feeding cats commercially available food, rather than raw, undercooked meat, prevents felines from becoming a host for oocysts, as higher prevalence is in areas where raw meat is fed. Researchers also suggest that owners restrict cats to live indoors and to be neutered or spayed to decrease stray cat populations and to reduce intermediate host interactions. It is suggested that fecal matter from litter boxes be collected daily, placed in a sealable bag, and disposed of in the trash rather than flushed in the toilet, so that water contamination is limited.
Studies have found that wetlands with a high density of vegetation decrease the concentration of oocysts in water through two possible mechanisms. Firstly, vegetation decreases flow velocities, which enables more settling because of increased transport time. Secondly, the vegetation can remove oocysts through its ability to mechanically strain the water, as well as through the process of adhesion (i.e. attachment to biofilms). Areas of erosion and destruction of coastal wetlands have been found to harbour increased concentrations of T. gondii oocysts, which then flow into open coastal waters. Current physical and chemical treatments typically utilized in water treatment facilities have been proven to be ineffective against T. gondii. Research has shown that UV-C disinfection of water containing oocysts results in inactivation and possible sterilization.
Genome
The genomes of more than 60 strains of T. gondii have been sequenced. Most are 60–80 Mb in size and consist of 11–14 chromosomes. The major strains encode 7,800–10,000 proteins, of which about 5,200 are conserved across RH, GT1, ME49, VEG. A database, ToxoDB, has been established to document genomic information on Toxoplasma.
History
In 1908, while working at the Pasteur Institute in Tunis, Charles Nicolle and Louis Manceaux discovered a protozoan organism in the tissues of a hamster-like rodent known as the gundi, Ctenodactylus gundi. Although Nicolle and Manceaux initially believed the organism to be a member of the genus Leishmania that they described as "Leishmania gondii", they soon realized they had discovered a new organism entirely; they renamed it Toxoplasma gondii. The new genus name Toxoplasma is a reference to its morphology: Toxo, from Greek (, 'arc, bow'), and (, 'shape, form') and the host in which it was discovered, the gundi (gondii). The same year Nicolle and Mancaeux discovered T. gondii, Alfonso Splendore identified the same organism in a rabbit in Brazil. However, he did not give it a name. In 1914, Italian tropicalist Aldo Castellani "was first to suspect that toxoplasmosis could affect humans".
The first conclusive identification of T. gondii in humans was in an infant girl delivered full term by Caesarean section on May 23, 1938, at Babies' Hospital in New York City. The girl began having seizures at three days of age, and doctors identified lesions in the maculae of both of her eyes. When she died at one month of age, an autopsy was performed. Lesions discovered in her brain and eye tissue were found to have both free and intracellular T. gondii. Infected tissue from the girl was homogenized and inoculated intracerebrally into rabbits and mice; they then developed encephalitis. Later, congenital transmission was confirmed in many other species, particularly infected sheep and rodents.
The possibility of T. gondii transmission via consumption of undercooked meat was first proposed by D. Weinman and A.H Chandler in 1954. In 1960, the relevant cyst wall were shown to dissolve in the proteolytic enzymes found in the stomach, releasing infectious bradyzoites into the stomach (which pass into the intestine). The hypothesis of transmission via consumption of undercooked meat was tested on children hospitalized in a sanatorium in Paris in 1965 by Georges Desmonts et al; incidence of T. gondii rose from 10% to 50% after a year of adding two portions of cooked-rare beef or horse meat to many children's daily diets, and to 100% among those fed cooked-rare lamb chops.
A 1959 Mumbai-based study found there prevalence in strict vegetarians was similar to that of non-vegetarians. This raised the possibility of a third major route of infection, beyond congenital and non well-cooked meat carnivorous transmission.
In 1970, oocysts were found in (cat) feces. The fecal–oral route of infection via oocysts was demonstrated. In the 1970s and 1980s feces of a vast range of infected animal species was tested to see if it contained oocysts—at least 17 species of felids shed oocysts, but no non-felid has been shown to allow T. gondii sexual reproduction (leading to oocyst shedding).
In 1984 Elmer R. Pfefferkorn published his discovery that treatment of human fibroblasts with human recombinant interferon gamma blocks the growth of T. gondii.
Behavioral differences of infected hosts
There are many instances where behavioural changes were reported in rodents with T. gondii. The changes seen were a reduction in their innate dislike of cats, which made it easier for cats to prey on the rodents. In an experiment conducted by Berdoy and colleagues, the infected rats showed preference for the cat odour area versus the area with the rabbit scent, therefore making it easier for the parasite to take its final step in its definitive feline host. This is an example of the extended phenotype concept, that is, the idea that the behaviour of the infected animal changes in order to maximize survival of the genes that increase predation of the intermediate rodent host.
Differences in sex-dependent behavior observed in infected hosts compared to non-infected individuals can be attributed to differences in testosterone. Infected males had higher levels of testosterone while infected females had significantly lower levels, compared to their non-infected equivalents.
Looking at humans, studies using the Cattell's 16 Personality Factor questionnaire found that infected men scored lower on Factor G (superego strength/rule consciousness) and higher on Factor L (vigilance) while the opposite pattern was observed for infected women. Such men were more likely to disregard rules and were more expedient, suspicious, and jealous. On the other hand, women were more warm-hearted, outgoing, conscientious, and moralistic.
Published research has also indicated that T. gondii infection could potentially promote changes in a person's political beliefs and values. Those who are infected with the parasite tend to exhibit a higher degree of "us versus them" thinking.
Mice infected with T. gondii have a worse motor performance than non-infected mice. Thus, a computerized simple reaction test was given to both infected and non-infected adults. It was found that the infected adults performed much more poorly and lost their concentration more quickly than the control group. But, the effect of the infection only explains less than 10% of the variability in performance (i.e., there could be other confounding factors).
Correlation has also been observed between seroprevalence of T. gondii in humans and increased risk of traffic accidents. Infected subjects have a 2.65 times higher risk of getting into a traffic accident. A Turkish study confirmed this holds true among drivers.
This parasite has been associated with many neurological disorders such as schizophrenia. In a meta-analysis of 23 studies that met inclusion criteria, the seroprevalence of antibodies to T. gondii in people with schizophrenia is significantly higher than in control populations (OR=2.73, P<0.000001).
A 2009 summary of studies found that suicide attempters had far more indicative (IgG) antibodies than mental health inpatients without a suicide attempt. Infection was also shown to be associated with suicide in women over the age of 60. (P<0.005)
Research on the linkage between T. gondii infection and entrepreneurial behavior showed that students who tested positive for T. gondii exposure were 1.4 times more likely to major in business and 1.7 times more likely to have an emphasis in "management and entrepreneurship". Among 197 participants of entrepreneurship events, T. gondii exposure was correlated with being 1.8 times more likely to have started their own business.
Another population-representative study with 7440 people in the United States found that Toxoplasma infection was 2.4 fold more common in people who had a history of manic and depression symptoms (bipolar disorder Type 1) compared to the general population.
As mentioned before, these results of increased proportions of people seropositive for the parasite in cases of these neurological disorders do not necessarily indicate a causal relationship between the infection and disorder. It is also important to mention that in 2016 a population-representative birth cohort study which was done, to test a hypothesis that toxoplasmosis is related to impairment in brain and behaviour measured by a range of phenotypes including neuropsychiatric disorders, poor impulse control, personality and neurocognitive deficits. The results of this study did not support the results in the previously mentioned studies, more than marginally. None of the P-values showed significance for any outcome measure. Thus, according to this study, the presence of T. gondii antibodies is not correlated to increase susceptibility to any of the behaviour phenotypes (except possibly to a higher rate of unsuccessful attempted suicide). This team did not observe any significant association between T. gondii seropositivity and schizophrenia. The team notes that the null findings might be a false negative due to low statistical power because of small sample sizes but against this weights that their setup should avoid some possibilities for errors in the about 40 studies that did show a positive correlation. They concluded that further studies should be performed.
The mechanism behind behavioral changes is partially attributed to increased dopamine metabolism, which can be neutralized by dopamine antagonist medications. T. gondii has two genes that code for a bifunctional phenylalanine and tyrosine hydroxylase, two important and rate-limiting steps of dopamine biosynthesis. One of the genes is constitutively expressed, while the other is only produced during cyst development. In addition to additional dopamine production, T. gondii infection also produces long-lasting epigenetic changes in animals that increase the expression of vasopressin, a probable cause of alterations that persist after the clearance of the infection.
In 2022, a study published in Communications Biology of a well-documented population of wolves studied throughout their lives, suggested that T. gondii also may have a significant effect on their behavior. It suggested that infection with this parasite emboldened infected wolves into behavior that determined leadership roles and influenced risk-taking behavior, perhaps even motivating establishment of new independent packs that they would establish and lead in behavior patterns differing from that of the packs into which they were born. The study determined that at times, an infected wolf would become the only breeding male in a pack, leading to a significant effect on another species by T. gondii.
Potential medical use
In July 2024, a study published in Nature Microbiology showed that T. gondii can be engineered to deliver the MECP2 protein, a therapeutic target of Rett syndrome, to the brain of infected mice.
| Biology and health sciences | SAR supergroup | Plants |
196741 | https://en.wikipedia.org/wiki/SARS | SARS | Severe acute respiratory syndrome (SARS) is a viral respiratory disease of zoonotic origin caused by the virus SARS-CoV-1, the first identified strain of the SARS-related coronavirus. The first known cases occurred in November 2002, and the syndrome caused the 2002–2004 SARS outbreak. In the 2010s, Chinese scientists traced the virus through the intermediary of Asian palm civets to cave-dwelling horseshoe bats in Xiyang Yi Ethnic Township, Yunnan.
SARS was a relatively rare disease; at the end of the epidemic in June 2003, the incidence was 8,422 cases with a case fatality rate (CFR) of 11%. No cases of SARS-CoV-1 have been reported worldwide since 2004.
In December 2019, a second strain of SARS-CoV was identified: SARS-CoV-2. This strain causes coronavirus disease 2019 (COVID-19), the disease behind the COVID-19 pandemic.
Signs and symptoms
SARS produces flu-like symptoms which may include fever, muscle pain, lethargy, cough, sore throat, and other nonspecific symptoms. SARS often leads to shortness of breath and pneumonia, which may be direct viral pneumonia or secondary bacterial pneumonia.
The average incubation period for SARS is four to six days, although it is rarely as short as one day or as long as 14 days.
Transmission
The primary route of transmission for SARS-CoV is contact of the mucous membranes with respiratory droplets or fomites. As with all respiratory pathogens once presumed to transmit via respiratory droplets, it is highly likely to be carried by the aerosols generated during routine breathing, talking, and even singing. While diarrhea is common in people with SARS, the fecal–oral route does not appear to be a common mode of transmission. The basic reproduction number of SARS-CoV, R0, ranges from 2 to 4 depending on different analyses. Control measures introduced in April 2003 reduced the R to 0.4.
Diagnosis
SARS-CoV may be suspected in a patient who has:
Any of the symptoms, including a fever of or higher, and
Either a history of:
Contact (sexual or casual) with someone with a diagnosis of SARS within the last 10 days or
Travel to any of the regions identified by the World Health Organization (WHO) as areas with recent local transmission of SARS.
Clinical criteria of Sars-CoV diagnosis
Early illness: equal to or more than 2 of the following: chills, rigors, myalgia, diarrhea, sore throat (self-reported or observed)
Mild-to-moderate illness: temperature of > plus indications of lower respiratory tract infection (cough, dyspnea)
Severe illness: ≥1 of radiographic evidence, presence of ARDS, autopsy findings in late patients.
For a case to be considered probable, a chest X-ray must be indicative for atypical pneumonia or acute respiratory distress syndrome.
The WHO has added the category of "laboratory confirmed SARS" which means patients who would otherwise be considered "probable" and have tested positive for SARS based on one of the approved tests (ELISA, immunofluorescence or PCR) but whose chest X-ray findings do not show SARS-CoV infection (e.g. ground glass opacities, patchy consolidations unilateral).
The appearance of SARS-CoV in chest X-rays is not always uniform but generally appears as an abnormality with patchy infiltrates.
Prevention
There is a vaccine for SARS, although in March 2020 immunologist Anthony Fauci said the CDC developed one and placed it in the Strategic National Stockpile. That vaccine is a final product and field-ready as of March 2022. Clinical isolation and vaccination remain the most effective means to prevent the spread of SARS. Other preventive measures include:
Hand-washing with soap and water, or use of alcohol-based hand sanitizer
Disinfection of surfaces of fomites to remove viruses
Avoiding contact with bodily fluids
Washing the personal items of someone with SARS in hot, soapy water (eating utensils, dishes, bedding, etc.)
Avoiding travel to affected areas
Wearing masks and gloves
Keeping people with symptoms home from school
Simple hygiene measures
Distancing oneself at least 6 feet if possible to minimize the chances of transmission of the virus
Many public health interventions were made to try to control the spread of the disease, which is mainly spread through respiratory droplets in the air, either inhaled or deposited on surfaces and subsequently transferred to a body's mucous membranes. These interventions included earlier detection of the disease; isolation of people who are infected; droplet and contact precautions; and the use of personal protective equipment (PPE), including masks and isolation gowns. A 2017 meta-analysis found that for medical professionals wearing N-95 masks could reduce the chances of getting sick up to 80% compared to no mask. A screening process was also put in place at airports to monitor air travel to and from affected countries.
SARS-CoV is most infectious in severely ill patients, which usually occurs during the second week of illness. This delayed infectious period meant that quarantine was highly effective; people who were isolated before day five of their illness rarely transmitted the disease to others.
As of 2017, the CDC was still working to make federal and local rapid-response guidelines and recommendations in the event of a reappearance of the virus.
Treatment
As SARS is a viral disease, antibiotics do not have direct effect but may be used against bacterial secondary infection. Treatment of SARS is mainly supportive with antipyretics, supplemental oxygen and mechanical ventilation as needed. While ribavirin is commonly used to treat SARS, there seems to have little to no effect on SARS-CoV, and no impact on patient's outcomes. There is currently no proven antiviral therapy. Tested substances, include ribavirin, lopinavir, ritonavir, type I interferon, that have thus far shown no conclusive contribution to the disease's course. Administration of corticosteroids, is recommended by the British Thoracic Society/British Infection Society/Health Protection Agency in patients with severe disease and O2 saturation of <90%.
People with SARS-CoV must be isolated, preferably in negative-pressure rooms, with complete barrier nursing precautions taken for any necessary contact with these patients, to limit the chances of medical personnel becoming infected. In certain cases, natural ventilation by opening doors and windows is documented to help decreasing indoor concentration of virus particles.
Some of the more serious damage caused by SARS may be due to the body's own immune system reacting in what is known as cytokine storm.
Vaccine
Vaccines can help the immune system to create enough antibodies and decrease a risk of side effects like arm pain, fever, and headache. According to research papers published in 2005 and 2006, the identification and development of novel vaccines and medicines to treat SARS was a priority for governments and public health agencies around the world. In early 2004, an early clinical trial on volunteers was planned. A major researcher's 2016 request, however, demonstrated that no field-ready SARS vaccine had been completed because likely market-driven priorities had ended funding.
Prognosis
Several consequent reports from China on some recovered SARS patients showed severe long-time sequelae. The most typical diseases include, among other things, pulmonary fibrosis, osteoporosis, and femoral necrosis, which have led in some cases to the complete loss of working ability or even self-care ability of people who have recovered from SARS. As a result of quarantine procedures, some of the post-SARS patients have been diagnosed with post-traumatic stress disorder (PTSD) and major depressive disorder.
Epidemiology
SARS was a relatively rare disease; at the end of the epidemic in June 2003, the incidence was 8,422 cases with a case fatality rate (CFR) of 11%.
The case fatality rate (CFR) ranges from 0% to 50% depending on the age group of the patient. Patients under 24 were least likely to die (less than 1%); those 65 and older were most likely to die (over 55%).
As with MERS and COVID-19, SARS resulted in significantly more deaths of males than females.
Outbreak in South China
The SARS epidemic began in the Guangdong province of China in November 2002. The earliest case developed symptoms on 16 November 2002. The index patient, a farmer from Shunde, Foshan, Guangdong, was treated in the First People's Hospital of Foshan. The patient died soon after, and no definite diagnosis was made on his cause of death. Despite taking some action to control it, Chinese government officials did not inform the World Health Organization of the outbreak until February 2003. This lack of openness caused delays in efforts to control the epidemic, resulting in criticism of the People's Republic of China from the international community. China officially apologized for early slowness in dealing with the SARS epidemic.
In 2003, when the virus broke out in China, a 72 year old with SARS infected multiple people on board an Air China Boeing 737, causing 5 deaths.
The viral outbreak was subsequently genetically traced to a colony of cave-dwelling horseshoe bats in Xiyang Yi Ethnic Township, Yunnan.
The outbreak first came to the attention of the international medical community on 27 November 2002, when Canada's Global Public Health Intelligence Network (GPHIN), an electronic warning system that is part of the World Health Organization's Global Outbreak Alert and Response Network (GOARN), picked up reports of a "flu outbreak" in China through Internet media monitoring and analysis and sent them to the WHO. While GPHIN's capability had recently been upgraded to enable Arabic, Chinese, English, French, Russian, and Spanish translation, the system was limited to English or French in presenting this information. Thus, while the first reports of an unusual outbreak were in Chinese, an English report was not generated until 21 January 2003. The first super-spreader was admitted to the Sun Yat-sen Memorial Hospital in Guangzhou on 31 January, which soon spread the disease to nearby hospitals.
In early April 2003, after a prominent physician, Jiang Yanyong, pushed to report the danger to China, there appeared to be a change in official policy when SARS began to receive a much greater prominence in the official media. Some have directly attributed this to the death of an American teacher, James Earl Salisbury, in Hong Kong. It was around this same time that Jiang Yanyong made accusations regarding the undercounting of cases in Beijing military hospitals. After intense pressure, Chinese officials allowed international officials to investigate the situation there. This revealed problems plaguing the aging mainland Chinese healthcare system, including increasing decentralization, red tape, and inadequate communication.
Many healthcare workers in the affected nations risked their lives and died by treating patients, and trying to contain the infection before ways to prevent infection were known.
Spread to other regions
The epidemic reached the public spotlight in February 2003, when an American businessman traveling from China, Johnny Chen, became affected by pneumonia-like symptoms while on a flight to Singapore. The plane stopped in Hanoi, Vietnam, where the patient died in Hanoi French Hospital. Several of the medical staff who treated him soon developed the same disease despite basic hospital procedures. Italian doctor Carlo Urbani identified the threat and communicated it to WHO and the Vietnamese government; he later died from the disease.
The severity of the symptoms and the infection among hospital staff alarmed global health authorities, who were fearful of another emergent pneumonia epidemic. On 12 March 2003, the WHO issued a global alert, followed by a health alert by the United States Centers for Disease Control and Prevention (CDC). Local transmission of SARS took place in Toronto, Ottawa, San Francisco, Ulaanbaatar, Manila, Singapore, Taiwan, Hanoi and Hong Kong whereas within China it spread to Guangdong, Jilin, Hebei, Hubei, Shaanxi, Jiangsu, Shanxi, Tianjin, and Inner Mongolia.
Hong Kong
The disease spread in Hong Kong from Liu Jianlun, a Guangdong doctor who was treating patients at Sun Yat-Sen Memorial Hospital. He arrived in February and stayed on the ninth floor of the Metropole Hotel in Kowloon, infecting 16 of the hotel visitors. Those visitors traveled to Canada, Singapore, Taiwan, and Vietnam, spreading SARS to those locations.
Another larger cluster of cases in Hong Kong centred on the Amoy Gardens housing estate. Its spread is suspected to have been facilitated by defects in its bathroom drainage system that allowed sewer gases including virus particles to vent into the room. Bathroom fans exhausted the gases and wind carried the contagion to adjacent downwind complexes. Concerned citizens in Hong Kong worried that information was not reaching people quickly enough and created a website called sosick.org, which eventually forced the Hong Kong government to provide information related to SARS in a timely manner. The first cohort of affected people were discharged from hospital on 29 March 2003.
Canada
The first case of SARS in Toronto was identified on 23 February 2003. Beginning with an elderly woman, Kwan Sui-Chu, who had returned from a trip to Hong Kong and died on 5 March, the virus eventually infected 257 individuals in the province of Ontario. The trajectory of this outbreak is typically divided into two phases, the first centring around her son Tse Chi Kwai, who infected other patients at the Scarborough Grace Hospital and died on 13 March. The second major wave of cases was clustered around accidental exposure among patients, visitors, and staff within the North York General Hospital. The WHO officially removed Toronto from its list of infected areas by the end of June 2003.
The official response by the Ontario provincial government and Canadian federal government has been widely criticized in the years following the outbreak. Brian Schwartz, vice-chair of Ontario's SARS Scientific Advisory Committee, described public health officials' preparedness and emergency response at the time of the outbreak as "very, very basic and minimal at best". Critics of the response often cite poorly outlined and enforced protocol for protecting healthcare workers and identifying infected patients as a major contributing factor to the continued spread of the virus. The atmosphere of fear and uncertainty surrounding the outbreak resulted in staffing issues in area hospitals when healthcare workers elected to resign rather than risk exposure to SARS.
Identification of virus
In late February 2003, Italian doctor Carlo Urbani was called into The French Hospital of Hanoi to look at Johnny Chen, an American businessman who had fallen ill with what doctors thought was a bad case of influenza. Urbani realized that Chen's ailment was probably a new and highly contagious disease. He immediately notified the WHO. He also persuaded the Vietnamese Health Ministry to begin isolating patients and screening travelers, thus slowing the early pace of the epidemic. He subsequently contracted the disease himself, and died in March 2003.
Malik Peiris and his colleagues became the first to isolate the virus that causes SARS, a novel coronavirus now known as SARS-CoV-1. By June 2003, Peiris, together with his long-time collaborators Leo Poon and Guan Yi, has developed a rapid diagnostic test for SARS-CoV-1 using real-time polymerase chain reaction. The CDC and Canada's National Microbiology Laboratory identified the SARS genome in April 2003. Scientists at Erasmus University in Rotterdam, the Netherlands demonstrated that the SARS coronavirus fulfilled Koch's postulates thereby suggesting it as the causative agent. In the experiments, macaques infected with the virus developed the same symptoms as human SARS patients.
Origin and animal vectors
In late May 2003, a study was conducted using samples of wild animals sold as food in the local market in Guangdong, China. The study found that "SARS-like" coronaviruses could be isolated from masked palm civets (Paguma sp.). Genomic sequencing determined that these animal viruses were very similar to human SARS viruses, however they were phylogenetically distinct, and so the study concluded that it was unclear whether they were the natural reservoir in the wild. Still, more than 10,000 masked palm civets were killed in Guangdong Province since they were a "potential infectious source." The virus was also later found in raccoon dogs (Nyctereuteus sp.), ferret badgers (Melogale spp.), and domestic cats.
In 2005, two studies identified a number of SARS-like coronaviruses in Chinese bats. Phylogenetic analysis of these viruses indicated a high probability that SARS coronavirus originated in bats and spread to humans either directly or through animals held in Chinese markets. The bats did not show any visible signs of disease, but are the likely natural reservoirs of SARS-like coronaviruses. In late 2006, scientists from the Chinese Centre for Disease Control and Prevention of Hong Kong University and the Guangzhou Centre for Disease Control and Prevention established a genetic link between the SARS coronavirus appearing in civets and in the second, 2004 human outbreak, bearing out claims that the disease had jumped across species.
It took 14 years to find the original bat population likely responsible for the SARS pandemic. In December 2017, "after years of searching across China, where the disease first emerged, researchers reported ... that they had found a remote cave in Xiyang Yi Ethnic Township, Yunnan province, which is home to horseshoe bats that carry a strain of a particular virus known as a coronavirus. This strain has all the genetic building blocks of the type that triggered the global outbreak of SARS in 2002." The research was performed by Shi Zhengli, Cui Jie, and co-workers at the Wuhan Institute of Virology, China, and published in PLOS Pathogens. The authors are quoted as stating that "another deadly outbreak of SARS could emerge at any time. The cave where they discovered their strain is only a kilometre from the nearest village." The virus was ephemeral and seasonal in bats. In 2019, a similar virus to SARS caused a cluster of infections in Wuhan, eventually leading to the COVID-19 pandemic.
A small number of cats and dogs tested positive for the virus during the outbreak. However, these animals did not transmit the virus to other animals of the same species or to humans.
Containment
The World Health Organization declared severe acute respiratory syndrome contained on 5 July 2003. The containment was achieved through successful public health measures. In the following months, four SARS cases were reported in China between December 2003 and January 2004.
While SARS-CoV-1 probably persists as a potential zoonotic threat in its original animal reservoir, human-to-human transmission of this virus may be considered eradicated because no human case has been documented since four minor, brief, subsequent outbreaks in 2004.
Laboratory accidents
After containment, there were four laboratory accidents that resulted in infections.
One postdoctoral student at the National University of Singapore in Singapore in August 2003
A 44-year-old senior scientist at the National Defense University in Taipei in December 2003. He was confirmed to have the SARS virus after working on a SARS study in Taiwan's only BSL-4 lab. The Taiwan CDC later stated the infection occurred due to laboratory misconduct.
Two researchers at the Chinese Institute of Virology in Beijing, China around April 2004, who spread it to around six other people. The two researchers contracted it 2 weeks apart.
Study of live SARS specimens requires a biosafety level 3 (BSL-3) facility; some studies of inactivated SARS specimens can be done at biosafety level 2 facilities.
Society and culture
Fear of contracting the virus from consuming infected wild animals resulted in public bans and reduced business for meat markets in southern China and Hong Kong. The WHO declared the end of the pandemic on 24 March 2004.
| Biology and health sciences | Infectious disease | null |
196782 | https://en.wikipedia.org/wiki/Kaffir%20lime | Kaffir lime | Citrus hystrix, called the kaffir lime, Thai lime or makrut lime, (, ) is a citrus fruit native to tropical Southeast Asia.
Its fruit and leaves are used in Southeast Asian cuisine, and its essential oil is used in perfumery. Its rind and crushed leaves emit an intense citrus fragrance.
Names
"Kaffir" is thought to ultimately derive from the Arabic kafir, meaning infidel, though the mechanism by which it came to be applied to the lime is uncertain. Following the takeover of the Swahili coast, Muslims used the term to refer to the non-Muslim indigenous Africans, who were increasingly abducted for the Indian Ocean slave trade, which reached a height in the fifteenth and sixteenth century.
The most likely etymology is through the Kaffirs, an ethnic group in Sri Lanka partly descended from enslaved Bantu. The earliest known reference, under the alternative spelling "caffre" is in the 1888 book The Cultivated Oranges, Lemons Etc. of India and Ceylon by Emanuel Bonavia, who notes, "The plantation coolies also smear it over their feet and legs, to keep off land leeches; and therefore in Ceylon [Sri Lanka] it has also got the name of Kudalu dchi, or Leech Lime. Europeans call it Caffre Lime." Similarly, H.F. MacMillan's 1910 book A Handbook of Tropical Gardening and Planting notes, "The 'Kaffir Lime' in Ceylon."
Another proposed etymology is directly by Indian Muslims of the imported fruit from the non-Muslim lands to the east to "convey otherness and exotic provenance." Claims that the name of the fruit derives directly from the South African ethnic slur "kaffir" (see "South Africa" below) are not well supported.
C. hystrix is known by various names in its native areas:
in Javanese and in Malay (respectively into Indonesian and Malaysian) both meaning "warty/rough-skinned lime" due to the fruit's bumpy texture.
( "arrow-leaf lime") in Chinese.
or in the Philippines. The city of Cabuyao in Laguna is named after the fruit.
Kolumichai, கொலுமிச்சை in Kongu Tamil
or (, ) in Thailand (a name also used for the bergamot orange).
(, ) in Laos.
or in Vietnam.
in Réunion Island
The micrantha, a similar citrus fruit native to the Philippines that is ancestral to several hybrid limes, such as the Key lime and Persian lime, may represent the same species as C. hystrix, but the genomic characterization of the kaffir lime has not been performed in sufficient detail to allow a definitive conclusion.
South Africa
In South Africa, the Arabic kafir was adopted by White colonialists as "kaffir," an ethnic slur for black African people. Consequently, some authors favour switching from "kaffir lime" to "makrut lime," a less well-known name, while in South Africa, it is usually referred to as "Thai lime".
Description
C. hystrix is a thorny shrub or small tree, tall, with aromatic and distinctively shaped "double" leaves. These hourglass-shaped leaves comprise the leaf blade plus a flattened, leaf-like stalk (botanically, a winged petiole). The fruit is rough and green and ripens to yellow; it is distinguished by its bumpy exterior and small size, approximately wide. The fruits have thick skins (pericarps) and taste very acidic and slightly bitter. Flowers can have four to five petals that are white in color and are fragrant.
History
Pierre Sonnerat (1748–1814) collected specimens of it in 1771–72, and it appears in Lamarck's Encyclopédie Méthodique (1796).
Makrut lime appears in texts under the name of kaffir lime in 1868, in Ceylon, where rubbing the juice onto legs and socks prevents leech bites. This could be a possible origin of the name leech lime.
Uses
Culinary
C. hystrix leaves are used in Southeast Asian cuisines such as Indonesian, Laotian, Cambodian, and Thai. The leaves are the most frequently used part of the plant, fresh, dried, or frozen. The leaves are widely used in Thai cuisine (for dishes such as tom yum) and Cambodian cuisine (for the base paste "krueng"). The leaves are used in Vietnamese cuisine to add fragrance to chicken dishes and to decrease the pungent odor when steaming snails. Also, in Vietnamese villages that harvest silkworms, the silkworms in the pupa stage are stir fried with the kaffir lime leaves. The leaves are used in Indonesian cuisine (especially Balinese cuisine and Javanese cuisine) for foods such as soto ayam and are used along with Indonesian bay leaf for chicken and fish. They are also found in Malaysian and Burmese cuisines.
The rind (peel) is commonly used in Lao and Thai curry paste, adding an aromatic, astringent flavor. The zest of the fruit, referred to as combava, is used in creole cuisine to impart flavor in infused rums and rougails in Mauritius, Réunion, and Madagascar. In Cambodia, the entire fruit is crystallized/candied for eating.
Medicinal
The juice and rinds of the peel are used in traditional medicine in some Asian countries; the fruit's juice is often used in shampoo and is believed to kill head lice.
Other uses
The juice is used as a cleanser for clothing and hair in Thailand and occasionally in Cambodia. Lustral water mixed with slices of the fruit is used in religious ceremonies in Cambodia.
Makrut lime oil is used as raw material in many fields, including pharmaceutical, agronomic, food, sanitary, cosmetic, and perfume industries. It is also used extensively in aromatherapy and as an essential ingredient in various cosmetic and beauty products.
Cultivation
C. hystrix is grown worldwide in suitable climates as a garden shrub for home fruit production. It is well suited to container gardens and for large garden pots on patios, terraces, and in conservatories.
Main constituents
The compound responsible for the characteristic aroma was identified as (–)-(S)-citronellal, which is contained in the leaf oil up to 80 percent; minor components include citronellol (10 percent), nerol and limonene.
From a stereochemical point of view, it is remarkable that makrut lime leaves contain only the (S) stereoisomer of citronellal, whereas its enantiomer, (+)-(R)-citronellal is found in both lemon balm and (to a lesser degree) lemon grass, (however, citronellal is only a trace component in the latter's essential oil).
Makrut lime fruit peel contains an essential oil comparable to lime fruit peel oil; its main components are limonene and β-pinene.
Toxicity
C. hystrix contains significant quantities of furanocoumarins, in both the peel and the pulp. Furanocoumarins are known to cause phytophotodermatitis, a potentially severe skin inflammation. Cases of phytophotodermatitis induced by external use of C. hystrix have been reported.
| Biology and health sciences | Citrus fruits | Plants |
196788 | https://en.wikipedia.org/wiki/Steam%20locomotive | Steam locomotive | A steam locomotive is a locomotive that provides the force to move itself and other vehicles by means of the expansion of steam. It is fuelled by burning combustible material (usually coal, oil or, rarely, wood) to heat water in the locomotive's boiler to the point where it becomes gaseous and its volume increases 1,700 times. Functionally, it is a steam engine on wheels.
In most locomotives, the steam is admitted alternately to each end of its cylinders in which pistons are mechanically connected to the locomotive's main wheels. Fuel and water supplies are usually carried with the locomotive, either on the locomotive itself or in a tender coupled to it. Variations in this general design include electrically powered boilers, turbines in place of pistons, and using steam generated externally.
Steam locomotives were first developed in the United Kingdom during the early 19th century and used for railway transport until the middle of the 20th century. Richard Trevithick built the first steam locomotive known to have hauled a load over a distance at Pen-y-darren in 1804, although he produced an earlier locomotive for trial at Coalbrookdale in 1802. Salamanca, built in 1812 by Matthew Murray for the Middleton Railway, was the first commercially successful steam locomotive. Locomotion No. 1, built by George Stephenson and his son Robert's company Robert Stephenson and Company, was the first steam locomotive to haul passengers on a public railway, the Stockton and Darlington Railway, in 1825. Rapid development ensued; in 1830 George Stephenson opened the first public inter-city railway, the Liverpool and Manchester Railway, after the success of Rocket at the 1829 Rainhill Trials had proved that steam locomotives could perform such duties. Robert Stephenson and Company was the pre-eminent builder of steam locomotives in the first decades of steam for railways in the United Kingdom, the United States, and much of Europe.
Towards the end of the steam era, a longstanding British emphasis on speed culminated in a record, still unbroken, of by LNER Class A4 4468 Mallard, however there are long-standing claims that the Pennsylvania Railroad class S1 achieved speeds upwards of 150 mph, though this was never officially proven. In the United States, larger loading gauges allowed the development of very large, heavy locomotives such as the Union Pacific Big Boy, which weighs and has a tractive effort of .
Beginning in the early 1900s, steam locomotives were gradually superseded by electric and diesel locomotives, with railways fully converting to electric and diesel power beginning in the late 1930s. The majority of steam locomotives were retired from regular service by the 1980s, although several continue to run on tourist and heritage lines.
History
Britain
The earliest railways employed horses to draw carts along rail tracks. In 1784, William Murdoch, a Scottish inventor, built a small-scale prototype of a steam road locomotive in Birmingham. A full-scale rail steam locomotive was proposed by William Reynolds around 1787. An early working model of a steam rail locomotive was designed and constructed by steamboat pioneer John Fitch in the US during 1794. Some sources claim Fitch's model was operable already by the 1780s and that he demonstrated his locomotive to George Washington. His steam locomotive used interior bladed wheels guided by rails or tracks. The model still exists at the Ohio Historical Society Museum in Columbus, US. The authenticity and date of this locomotive is disputed by some experts and a workable steam train would have to await the invention of the high-pressure steam engine by Richard Trevithick, who pioneered the use of steam locomotives.
The first full-scale working railway steam locomotive was the gauge Coalbrookdale Locomotive built by Trevithick in 1802. It was constructed for the Coalbrookdale ironworks in Shropshire in the United Kingdom though no record of it working there has survived. On 21 February 1804, the first recorded steam-hauled railway journey took place as another of Trevithick's locomotives hauled a train along the -wide tramway from the Pen-y-darren ironworks, near Merthyr Tydfil, to Abercynon in South Wales. Accompanied by Andrew Vivian, it ran with mixed success. The design incorporated a number of important innovations that included using high-pressure steam which reduced the weight of the engine and increased its efficiency.
Trevithick visited the Newcastle area in 1804 and had a ready audience of colliery (coal mine) owners and engineers. The visit was so successful that the colliery railways in north-east England became the leading centre for experimentation and development of the steam locomotive. Trevithick continued his own steam propulsion experiments through another trio of locomotives, concluding with the Catch Me Who Can in 1808, first in the world to haul fare-paying passengers.
In 1812, Matthew Murray's successful twin-cylinder rack locomotive Salamanca first ran on the edge-railed rack-and-pinion Middleton Railway. Another well-known early locomotive was Puffing Billy, built 1813–14 by engineer William Hedley. It was intended to work on the Wylam Colliery near Newcastle upon Tyne. This locomotive is the oldest preserved, and is on static display at the Science Museum, London.
George Stephenson
George Stephenson, a former miner working as an engine-wright at Killingworth Colliery, developed up to sixteen Killingworth locomotives, including Blücher in 1814, another in 1815, and a (newly identified) Killingworth Billy in 1816. He also constructed The Duke in 1817 for the Kilmarnock and Troon Railway, which was the first steam locomotive to work in Scotland.
In 1825, Stephenson built Locomotion No. 1 for the Stockton and Darlington Railway, north-east England, which was the first public steam railway in the world. In 1829, his son Robert built in Newcastle The Rocket, which was entered in and won the Rainhill Trials. This success led to the company emerging as the pre-eminent builder of steam locomotives used on railways in the UK, US and much of Europe. The Liverpool and Manchester Railway opened a year later making exclusive use of steam power for passenger and goods trains.
United States
Before the arrival of British imports, some domestic steam locomotive prototypes were built and tested in the United States, including John Fitch's miniature prototype. A prominent full sized example was Col. John Steven's "steam wagon" which was demonstrated on a loop of track in Hoboken, New Jersey in 1825.
Many of the earliest locomotives for commercial use on American railroads were imported from Great Britain, including first the Stourbridge Lion and later the John Bull. However, a domestic locomotive-manufacturing industry was soon established. In 1830, the Baltimore and Ohio Railroad's Tom Thumb, designed by Peter Cooper, was the first commercial US-built locomotive to run in America; it was intended as a demonstration of the potential of steam traction rather than as a revenue-earning locomotive. The DeWitt Clinton, built in 1831 for the Mohawk and Hudson Railroad, was a notable early locomotive.
, the original John Bull was on static display in the National Museum of American History in Washington, D.C. The replica is preserved at the Railroad Museum of Pennsylvania.
Continental Europe
The first railway service outside the United Kingdom and North America was opened in 1829 in France between Saint-Etienne and Lyon; it was initially limited to animal traction and converted to steam traction early 1831, using Seguin locomotives. The first steam locomotive in service in Europe outside of France was named The Elephant, which on 5 May 1835 hauled a train on the first line in Belgium, linking Mechelen and Brussels.
In Germany, the first working steam locomotive was a rack-and-pinion engine, similar to the Salamanca, designed by the British locomotive pioneer John Blenkinsop. Built in June 1816 by Johann Friedrich Krigar in the Royal Berlin Iron Foundry (Königliche Eisengießerei zu Berlin), the locomotive ran on a circular track in the factory yard. It was the first locomotive to be built on the European mainland and the first steam-powered passenger service; curious onlookers could ride in the attached coaches for a fee. It is portrayed on a New Year's badge for the Royal Foundry dated 1816. Another locomotive was built using the same system in 1817. They were to be used on pit railways in Königshütte and in Luisenthal on the Saar (today part of Völklingen), but neither could be returned to working order after being dismantled, moved and reassembled. On 7 December 1835, the Adler ran for the first time between Nuremberg and Fürth on the Bavarian Ludwig Railway. It was the 118th engine from the locomotive works of Robert Stephenson and stood under patent protection.
In Russia, the first steam locomotive was built in 1834 by Cherepanovs, however, it suffered from the lack of coal in the area and was replaced with horse traction after all the woods nearby had been cut down. The first Russian Tsarskoye Selo steam railway started in 1837 with locomotives purchased from Robert Stephenson and Company.
In 1837, the first steam railway started in Austria on the Emperor Ferdinand Northern Railway between Vienna-Floridsdorf and Deutsch-Wagram. The oldest continually working steam engine in the world also runs in Austria: the GKB 671 built in 1860, has never been taken out of service, and is still used for special excursions.
In 1838, the third steam locomotive to be built in Germany, the Saxonia, was manufactured by the Maschinenbaufirma Übigau near Dresden, built by Prof. Johann Andreas Schubert. The first independently designed locomotive in Germany was the Beuth, built by August Borsig in 1841. The first locomotive produced by Henschel-Werke in Kassel, the Drache, was delivered in 1848.
The first steam locomotives operating in Italy were the Bayard and the Vesuvio, running on the Napoli-Portici line, in the Kingdom of the Two Sicilies.
The first railway line over Swiss territory was the Strasbourg–Basel line opened in 1844. Three years later, in 1847, the first fully Swiss railway line, the Spanisch Brötli Bahn, from Zürich to Baden was opened.
Australia
The arid nature of south Australia posed distinctive challenges to their early steam locomotion network. The high concentration of magnesium chloride in the well water (bore water) used in locomotive boilers on the Trans-Australian Railway caused serious and expensive maintenance problems. At no point along its route does the line cross a permanent freshwater watercourse, so bore water had to be relied on. No inexpensive treatment for the highly mineralised water was available, and locomotive boilers were lasting less than a quarter of the time normally expected. In the days of steam locomotion, about half the total train load was water for the engine. The line's operator, Commonwealth Railways, was an early adopter of the diesel-electric locomotive.
Components
Boiler
The fire-tube boiler was standard practice for steam locomotive. Although other types of boiler were evaluated they were not widely used, except for some 1,000 locomotives in Hungary which used the water-tube Brotan boiler.
A boiler consists of a firebox where the fuel is burned, a barrel where water is turned into steam, and a smokebox which is kept at a slightly lower pressure than outside the firebox.
Solid fuel, such as wood, coal or coke, is thrown into the firebox through a door by a fireman, onto a set of grates which hold the fuel in a bed as it burns. Ash falls through the grate into an ashpan. If oil is used as the fuel, a door is needed for adjusting the air flow, maintaining the firebox, and cleaning the oil jets.
The fire-tube boiler has internal tubes connecting the firebox to the smokebox through which the combustion gases flow transferring heat to the water. All the tubes together provide a large contact area, called the tube heating surface, between the gas and water in the boiler. Boiler water surrounds the firebox to stop the metal from becoming too hot. This is another area where the gas transfers heat to the water and is called the firebox heating surface. Ash and char collect in the smokebox as the gas gets drawn up the chimney (stack or smokestack in the US) by the exhaust steam from the cylinders.
The pressure in the boiler has to be monitored using a gauge mounted in the cab. Steam pressure can be released manually by the driver or fireman. If the pressure reaches the boiler's design working limit, a safety valve opens automatically to reduce the pressure and avoid a catastrophic accident.
The exhaust steam from the engine cylinders shoots out of a nozzle pointing up the chimney in the smokebox. The steam entrains or drags the smokebox gases with it which maintains a lower pressure in the smokebox than that under the firebox grate. This pressure difference causes air to flow up through the coal bed and keeps the fire burning.
The search for thermal efficiency greater than that of a typical fire-tube boiler led engineers, such as Nigel Gresley, to consider the water-tube boiler. Although he tested the concept on the LNER Class W1, the difficulties during development exceeded the will to increase efficiency by that route.
The steam generated in the boiler not only moves the locomotive, but is also used to operate other devices such as the whistle, the air compressor for the brakes, the pump for replenishing the water in the boiler and the passenger car heating system. The constant demand for steam requires a periodic replacement of water in the boiler. The water is kept in a tank in the locomotive tender or wrapped around the boiler in the case of a tank locomotive. Periodic stops are required to refill the tanks; an alternative was a scoop installed under the tender that collected water as the train passed over a track pan located between the rails.
While the locomotive is producing steam, the amount of water in the boiler is constantly monitored by looking at the water level in a transparent tube, or sight glass. Efficient and safe operation of the boiler requires keeping the level in between lines marked on the sight glass. If the water level is too high, steam production falls, efficiency is lost and water is carried out with the steam into the cylinders, possibly causing mechanical damage. More seriously, if the water level gets too low, the crown sheet (top sheet) of the firebox becomes exposed. Without water on top of the sheet to transfer away the heat of combustion, it softens and fails, letting high-pressure steam into the firebox and the cab. The development of the fusible plug, a temperature-sensitive device, ensured a controlled venting of steam into the firebox to warn the fireman to add water.
Scale builds up in the boiler and prevents adequate heat transfer, and corrosion eventually degrades the boiler materials to the point where it needs to be rebuilt or replaced. Start-up on a large engine may take hours of preliminary heating of the boiler water before sufficient steam is available.
Although the boiler is typically placed horizontally, for locomotives designed to work in locations with steep slopes it may be more appropriate to consider a vertical boiler or one mounted such that the boiler remains horizontal but the wheels are inclined to suit the slope of the rails.
Steam circuit
The steam generated in the boiler fills the space above the water in the partially filled boiler. Its maximum working pressure is limited by spring-loaded safety valves. It is then collected either in a perforated tube fitted above the water level or by a dome that often houses the regulator valve, or throttle, the purpose of which is to control the amount of steam leaving the boiler. The steam then either travels directly along and down a steam pipe to the engine unit or may first pass into the wet header of a superheater, the role of the latter being to improve thermal efficiency and eliminate water droplets suspended in the "saturated steam", the state in which it leaves the boiler. On leaving the superheater, the steam exits the dry header of the superheater and passes down a steam pipe, entering the steam chests adjacent to the cylinders of a reciprocating engine. Inside each steam chest is a sliding valve that distributes the steam via ports that connect the steam chest to the ends of the cylinder space. The role of the valves is twofold: admission of each fresh dose of steam, and exhaust of the used steam once it has done its work.
The cylinders are double-acting, with steam admitted to each side of the piston in turn. In a two-cylinder locomotive, one cylinder is located on each side of the vehicle. The cranks are set 90° out of phase. During a full rotation of the driving wheel, steam provides four power strokes; each cylinder receives two injections of steam per revolution. The first stroke is to the front of the piston and the second stroke to the rear of the piston; hence two working strokes. Consequently, two deliveries of steam onto each piston face in the two cylinders generates a full revolution of the driving wheel. Each piston is attached to the driving axle on each side by a connecting rod, and the driving wheels are connected together by coupling rods to transmit power from the main driver to the other wheels. Note that at the two "dead centres", when the connecting rod is on the same axis as the crankpin on the driving wheel, the connecting rod applies no torque to the wheel. Therefore, if both cranksets could be at "dead centre" at the same time, and the wheels should happen to stop in this position, the locomotive could not start moving. Therefore, the crankpins are attached to the wheels at a 90° angle to each other, so only one side can be at dead centre at a time.
Each piston transmits power through a crosshead, connecting rod (Main rod in the US) and a crankpin on the driving wheel (Main driver in the US) or to a crank on a driving axle. The movement of the valves in the steam chest is controlled through a set of rods and linkages called the valve gear, actuated from the driving axle or from the crankpin; the valve gear includes devices that allow reversing the engine, adjusting valve travel and the timing of the admission and exhaust events. The cut-off point determines the moment when the valve blocks a steam port, "cutting off" admission steam and thus determining the proportion of the stroke during which steam is admitted into the cylinder; for example a 50% cut-off admits steam for half the stroke of the piston. The remainder of the stroke is driven by the expansive force of the steam. Careful use of cut-off provides economical use of steam and in turn, reduces fuel and water consumption. The reversing lever (Johnson bar in the US), or screw-reverser (if so equipped), that controls the cut-off, therefore, performs a similar function to a gearshift in an automobile – maximum cut-off, providing maximum tractive effort at the expense of efficiency, is used to pull away from a standing start, whilst a cut-off as low as 10% is used when cruising, providing reduced tractive effort, and therefore lower fuel/water consumption.
Exhaust steam is directed upwards out of the locomotive through the chimney, by way of a nozzle called a blastpipe, creating the familiar "chuffing" sound of the steam locomotive. The blastpipe is placed at a strategic point inside the smokebox that is at the same time traversed by the combustion gases drawn through the boiler and grate by the action of the steam blast. The combining of the two streams, steam and exhaust gases, is crucial to the efficiency of any steam locomotive, and the internal profiles of the chimney (or, strictly speaking, the ejector) require careful design and adjustment. This has been the object of intensive studies by a number of engineers (and often ignored by others, sometimes with catastrophic consequences). The fact that the draught depends on the exhaust pressure means that power delivery and power generation are automatically self-adjusting. Among other things, a balance has to be struck between obtaining sufficient draught for combustion whilst giving the exhaust gases and particles sufficient time to be consumed. In the past, a strong draught could lift the fire off the grate, or cause the ejection of unburnt particles of fuel, dirt and pollution for which steam locomotives had an unenviable reputation. Moreover, the pumping action of the exhaust has the counter-effect of exerting back pressure on the side of the piston receiving steam, thus slightly reducing cylinder power. Designing the exhaust ejector became a specific science, with engineers such as Chapelon, Giesl and Porta making large improvements in thermal efficiency and a significant reduction in maintenance time and pollution. A similar system was used by some early gasoline/kerosene tractor manufacturers (Advance-Rumely/Hart-Parr) – the exhaust gas volume was vented through a cooling tower, allowing the steam exhaust to draw more air past the radiator.
Running gear
Running gear includes the brake gear, wheel sets, axleboxes, springing and the motion that includes connecting rods and valve gear. The transmission of the power from the pistons to the rails and the behaviour of the locomotive as a vehicle, being able to negotiate curves, points and irregularities in the track, is of paramount importance. Because reciprocating power has to be directly applied to the rail from 0 rpm upwards, this creates the problem of adhesion of the driving wheels to the smooth rail surface. Adhesive weight is the portion of the locomotive's weight bearing on the driving wheels. This is made more effective if a pair of driving wheels is able to make the most of its axle load, i.e. its individual share of the adhesive weight. Equalising beams connecting the ends of leaf springs have often been deemed a complication in Britain, however, locomotives fitted with the beams have usually been less prone to loss of traction due to wheel-slip. Suspension using equalizing levers between driving axles, and between driving axles and trucks, was standard practice on North American locomotives to maintain even wheel loads when operating on uneven track.
Locomotives with total adhesion, where all of the wheels are coupled together, generally lack stability at speed. To counter this, locomotives often fit unpowered carrying wheels mounted on two-wheeled trucks or four-wheeled bogies centred by springs/inverted rockers/geared rollers that help to guide the locomotive through curves. These usually take on weight – of the cylinders at the front or the firebox at the rear – when the width exceeds that of the mainframes. Locomotives with multiple coupled-wheels on a rigid chassis would have unacceptable flange forces on tight curves giving excessive flange and rail wear, track spreading and wheel climb derailments. One solution was to remove or thin the flanges on an axle. More common was to give axles end-play and use lateral motion control with spring or inclined-plane gravity devices.
Railroads generally preferred locomotives with fewer axles, to reduce maintenance costs. The number of axles required was dictated by the maximum axle loading of the railroad in question. A builder would typically add axles until the maximum weight on any one axle was acceptable to the railroad's maximum axle loading. A locomotive with a wheel arrangement of two lead axles, two drive axles, and one trailing axle was a high-speed machine. Two lead axles were necessary to have good tracking at high speeds. Two drive axles had a lower reciprocating mass than three, four, five or six coupled axles. They were thus able to turn at very high speeds due to the lower reciprocating mass. A trailing axle was able to support a huge firebox, hence most locomotives with the wheel arrangement of (American Type Atlantic) were called free steamers and were able to maintain steam pressure regardless of throttle setting.
Chassis
The chassis, or locomotive frame, is the principal structure onto which the boiler is mounted and which incorporates the various elements of the running gear. The boiler is rigidly mounted on a "saddle" beneath the smokebox and in front of the boiler barrel, but the firebox at the rear is allowed to slide forward and backwards, to allow for expansion when hot.
European locomotives usually use "plate frames", where two vertical flat plates form the main chassis, with a variety of spacers and a buffer beam at each end to form a rigid structure. When inside cylinders are mounted between the frames, the plate frames are a single large casting that forms a major support element. The axleboxes slide up and down to give some sprung suspension, against thickened webs attached to the frame, called "hornblocks".
American practice for many years was to use built-up bar frames, with the smokebox saddle/cylinder structure and drag beam integrated therein. In the 1920s, with the introduction of "superpower", the cast-steel locomotive bed became the norm, incorporating frames, spring hangers, motion brackets, smokebox saddle and cylinder blocks into a single complex, sturdy but heavy casting. A SNCF design study using welded tubular frames gave a rigid frame with a 30% weight reduction.
Fuel and water
Generally, the largest locomotives are permanently coupled to a tender that carries the water and fuel. Often, locomotives working shorter distances do not have a tender and carry the fuel in a bunker, with the water carried in tanks placed next to the boiler. The tanks can be in various configurations, including two tanks alongside (side tanks or pannier tanks), one on top (saddle tank) or one between the frames (well tank).
The fuel used depended on what was economically available to the railway. In the UK and other parts of Europe, plentiful supplies of coal made this the obvious choice from the earliest days of the steam engine. Until 1870, the majority of locomotives in the United States burned wood, but as the Eastern forests were cleared, coal gradually became more widely used until it became the dominant fuel worldwide in steam locomotives. Railways serving sugar cane farming operations burned bagasse, a byproduct of sugar refining. In the US, the ready availability and low price of oil made it a popular steam locomotive fuel after 1900 for the southwestern railroads, particularly the Southern Pacific. In the Australian state of Victoria, many steam locomotives were converted to heavy oil firing after World War II. German, Russian, Australian and British railways experimented with using coal dust to fire locomotives.
During World War 2, a number of Swiss steam shunting locomotives were modified to use electrically heated boilers, consuming around 480 kW of power collected from an overhead line with a pantograph. These locomotives were significantly less efficient than electric ones; they were used because Switzerland was suffering a coal shortage because of the War, but had access to plentiful hydroelectricity.
A number of tourist lines and heritage locomotives in Switzerland, Argentina and Australia have used light diesel-type oil.
Water was supplied at stopping places and locomotive depots from a dedicated water tower connected to water cranes or gantries. In the UK, the US and France, water troughs (track pans in the US) were provided on some main lines to allow locomotives to replenish their water supply without stopping, from rainwater or snowmelt that filled the trough due to inclement weather. This was achieved by using a deployable "water scoop" fitted under the tender or the rear water tank in the case of a large tank engine; the fireman remotely lowered the scoop into the trough, the speed of the engine forced the water up into the tank, and the scoop was raised again once it was full.
Water is essential for the operation of a steam locomotive. As Swengel argued:
Swengel went on to note that "at low temperature and relatively low boiler outputs", good water and regular boiler washout was an acceptable practice, even though such maintenance was high. As steam pressures increased, however, a problem of "foaming" or "priming" developed in the boiler, wherein dissolved solids in the water formed "tough-skinned bubbles" inside the boiler, which in turn were carried into the steam pipes and could blow off the cylinder heads. To overcome the problem, hot mineral-concentrated water was deliberately wasted (blown down) from the boiler periodically. Higher steam pressures required more blowing-down of water out of the boiler. Oxygen generated by boiling water attacks the boiler, and with increased steam pressure the rate of rust (iron oxide) generated inside the boiler increases. One way to help overcome the problem was water treatment. Swengel suggested that these problems contributed to the interest in electrification of railways.
In the 1970s, L.D. Porta developed a sophisticated system of heavy-duty chemical water treatment (Porta Treatment) that not only keeps the inside of the boiler clean and prevents corrosion, but modifies the foam in such a way as to form a compact "blanket" on the water surface that filters the steam as it is produced, keeping it pure and preventing carry-over into the cylinders of water and suspended abrasive matter.
Some steam locomotives have been run on alternative fuels such as used cooking oil like Grand Canyon Railway 4960, Grand Canyon Railway 29, U.S. Sugar 148, and the Disneyland Railroad Locomotives.
Crew
A steam locomotive is normally controlled from the boiler's backhead, and the crew is usually protected from the elements by a cab. A crew of at least two people is normally required to operate a steam locomotive. One, the train driver or engineer (North America), is responsible for controlling the locomotive's starting, stopping, and speed, and the fireman is responsible for maintaining the fire, regulating steam pressure and monitoring boiler and tender water levels. Due to the historical loss of operational infrastructure and staffing, preserved steam locomotives operating on the mainline will often have a support crew travelling with the train.
Fittings and appliances
All locomotives are fitted with a variety of appliances. Some of these relate directly to the operation of the steam engine; others are for signalling, train control or other purposes. In the United States, the Federal Railroad Administration mandated the use of certain appliances over the years in response to safety concerns. The most typical appliances are as follows:
Steam pumps and injectors
Water (feedwater) must be delivered to the boiler to replace that which is exhausted as steam after delivering a working stroke to the pistons. As the boiler is under pressure during operation, feedwater must be forced into the boiler at a pressure that is greater than the steam pressure, necessitating the use of some sort of pump. Hand-operated pumps sufficed for the very earliest locomotives. Later engines used pumps driven by the motion of the pistons (axle pumps), which were simple to operate, reliable and could handle large quantities of water but only operated when the locomotive was moving and could overload the valve gear and piston rods at high speeds. Steam injectors later replaced the pump, while some engines transitioned to turbopumps. Standard practice evolved to use two independent systems for feeding water to the boiler; either two steam injectors or, on more conservative designs, axle pumps when running at service speed and a steam injector for filling the boiler when stationary or at low speeds. By the 20th century virtually all new-built locomotives used only steam injectors – often one injector was supplied with "live" steam straight from the boiler itself and the other used exhaust steam from the locomotive's cylinders, which was more efficient (since it made use of otherwise wasted steam) but could only be used when the locomotive was in motion and the regulator was open. Injectors became unreliable if the feedwater was at a high temperature, so locomotives with feedwater heaters, tank locomotives with the tanks in contact with the boiler and condensing locomotives sometimes used reciprocating steam pumps or turbopumps.
Vertical glass tubes, known as water gauges or water glasses, show the level of water in the boiler and are carefully monitored at all times while the boiler is being fired. Before the 1870s it was more common to have a series of try-cocks fitted to the boiler within reach of the crew; each try cock (at least two and usually three were fitted) was mounted at a different level. By opening each try-cock and seeing if steam or water vented through it, the level of water in the boiler could be estimated with limited accuracy. As boiler pressures increased the use of try-cocks became increasingly dangerous and the valves were prone to blockage with scale or sediment, giving false readings. This led to their replacement with the sight glass. As with the injectors, two glasses with separate fittings were usually installed to provide independent readings.
Boiler insulation
The term for pipe and boiler insulation is "lagging" which derives from the cooper's term for a wooden barrel stave. Two of the earliest steam locomotives used wooden lagging to insulate their boilers: the Salamanca, the first commercially successful steam locomotive, built in 1812, and the Locomotion No. 1, the first steam locomotive to carry passengers on a public rail line. Large amounts of heat are wasted if a boiler is not insulated. Early locomotives used lags, shaped wooden staves, fitted lengthways along the boiler barrel, and held in place by hoops, metal bands, the terms and methods are from cooperage.
Improved insulating methods included applying a thick paste containing a porous mineral such as kieselgur, or attaching shaped blocks of insulating compound such as magnesia blocks. In the latter days of steam, "mattresses" of stitched asbestos cloth stuffed with asbestos fibre were fixed to the boiler, on separators so as not quite to touch the boiler. However, asbestos is currently banned in most countries for health reasons. The most common modern-day material is glass wool, or wrappings of aluminium foil.
The lagging is protected by a close-fitted sheet-metal casing known as boiler clothing or cleading.
Effective lagging is particularly important for fireless locomotives; however, in recent times under the influence of L.D. Porta, "exaggerated" insulation has been practised for all types of locomotive on all surfaces liable to dissipate heat, such as cylinder ends and facings between the cylinders and the mainframes. This considerably reduces engine warmup time with a marked increase in overall efficiency.
Safety valves
Early locomotives were fitted with a valve controlled by a weight suspended from the end of a lever, with the steam outlet being stopped by a cone-shaped valve. As there was nothing to prevent the weighted lever from bouncing when the locomotive ran over irregularities in the track, thus wasting steam, the weight was later replaced by a more stable spring-loaded column, often supplied by Salter, a well-known spring scale manufacturer. The danger of these devices was that the driving crew could be tempted to add weight to the arm to increase pressure. Most early boilers were fitted with a tamper-proof "lockup" direct-loaded ball valve protected by a cowl. In the late 1850s, John Ramsbottom introduced a safety valve that became popular in Britain during the latter part of the 19th century. Not only was this valve tamper-proof, but tampering by the driver could only have the effect of easing pressure. George Richardson's safety valve was an American invention introduced in 1875, and was designed to release the steam only at the moment when the pressure attained the maximum permitted. This type of valve is in almost universal use at present. Britain's Great Western Railway was a notable exception to this rule, retaining the direct-loaded type until the end of its separate existence, because it was considered that such a valve lost less pressure between opening and closing.
Pressure gauge
The earliest locomotives did not show the pressure of steam in the boiler, but it was possible to estimate this by the position of the safety valve arm which often extended onto the firebox back plate; gradations marked on the spring column gave a rough indication of the actual pressure. The promoters of the Rainhill trials urged that each contender have a proper mechanism for reading the boiler pressure, and Stephenson devised a nine-foot vertical tube of mercury with a sight-glass at the top, mounted alongside the chimney, for his Rocket. The Bourdon tube gauge, in which the pressure straightens an oval-section coiled tube of brass or bronze connected to a pointer, was introduced in 1849 and quickly gained acceptance, and is still used today. Some locomotives have an additional pressure gauge in the steam chest. This helps the driver avoid wheel-slip at startup, by warning if the regulator opening is too great.
Spark arrestors and smokeboxes
Spark arrestor and self-cleaning smokebox
Wood-burners emit large quantities of flying sparks which necessitate an efficient spark-arresting device generally housed in the smokestack. Many different types were fitted, the most common early type being the Bonnet stack that incorporated a cone-shaped deflector placed before the mouth of the chimney pipe, and a wire screen covering the wide stack exit. A more efficient design was the Radley and Hunter centrifugal stack patented in 1850 (commonly known as the diamond stack), incorporating baffles so oriented as to induce a swirl effect in the chamber that encouraged the embers to burn out and fall to the bottom as ash. In the self-cleaning smokebox the opposite effect was achieved: by allowing the flue gasses to strike a series of deflector plates, angled in such a way that the blast was not impaired, the larger particles were broken into small pieces that would be ejected with the blast, rather than settle in the bottom of the smokebox to be removed by hand at the end of the run. As with the arrestor, a screen was incorporated to retain any large embers.
Locomotives of the British Railways standard classes fitted with self-cleaning smokeboxes were identified by a small cast oval plate marked "S.C.", fitted at the bottom of the smokebox door. These engines required different disposal procedures and the plate highlighted this need to depot staff.
Stokers
A factor that limits locomotive performance is the rate at which fuel is fed into the fire. In the early 20th century some locomotives became so large that the fireman could not shovel coal fast enough. In the United States, various steam-powered mechanical stokers became standard equipment and were adopted and used elsewhere including Australia and South Africa.
Feedwater heating
Introducing cold water into a boiler reduces power, and from the 1920s a variety of heaters were incorporated. The most common type for locomotives was the exhaust steam feedwater heater that piped some of the exhaust through small tanks mounted on top of the boiler or smokebox or into the tender tank; the warm water then had to be delivered to the boiler by a small auxiliary steam pump. The rare economiser type differed in that it extracted residual heat from the exhaust gases. An example of this is the pre-heater drum(s) found on the Franco-Crosti boiler.
The use of live steam and exhaust steam injectors also assists in the pre-heating of boiler feedwater to a small degree, though there is no efficiency advantage to live steam injectors. Such pre-heating also reduces the thermal shock that a boiler might experience when cold water is introduced directly. This is further helped by the top feed, where water is introduced to the highest part of the boiler and made to trickle over a series of trays. George Jackson Churchward fitted this arrangement to the high end of his domeless coned boilers. Other British lines such as the London, Brighton & South Coast Railway fitted some locomotives with the top feed inside a separate dome forward of the main one.
Condensers and water re-supply
Steam locomotives consume vast quantities of water because they operate on an open cycle, expelling their steam immediately after a single use rather than recycling it in a closed loop as stationary and marine steam engines do. Water was a constant logistical problem, and condensing engines were devised for use in desert areas. These engines had huge radiators in their tenders and instead of exhausting steam out of the funnel it was captured, passed back to the tender and condensed. The cylinder lubricating oil was removed from the exhausted steam to avoid a phenomenon known as priming, a condition caused by foaming in the boiler which would allow water to be carried into the cylinders causing damage because of its incompressibility. The most notable engines employing condensers (Class 25, the "puffers which never puff") worked across the Karoo desert of South Africa from the 1950s until the 1980s.
Some British and American locomotives were equipped with scoops which collected water from "water troughs" (track pans in the US) while in motion, thus avoiding stops for water. In the US, small communities often did not have refilling facilities. During the early days of railroading, the crew simply stopped next to a stream and filled the tender using leather buckets. This was known as "jerking water" and led to the term "jerkwater towns" (meaning a small town, a term which today is considered derisive). In Australia and South Africa, locomotives in drier regions operated with large oversized tenders and some even had an additional water wagon, sometimes called a "canteen" or in Australia (particularly in New South Wales) a "water gin".
Steam locomotives working on underground railways (such as London's Metropolitan Railway) were fitted with condensing apparatus to prevent steam from escaping into the railway tunnels. These were still being used between King's Cross and Moorgate into the early 1960s.
Braking
Steam locomotives usually have their own braking system, independent from the rest of the train. Locomotive brakes employ large shoes which press against the driving wheel treads. They can be either air brakes or steam brakes. In addition, they nearly always have a handbrake to keep the locomotive stationary when there is no steam pressure to power the other braking systems.
Because of the limited braking force provided by locomotive-only brakes, many steam locomotives were fitted with a train brake. These came in two main varieties; air brakes and vacuum brakes. These allowed the driver to control the brakes on all cars in the train.
Air brakes, invented by George Westinghouse, use a steam-driven air compressor mounted on the side of the boiler to create the compressed air needed to power the brake system. Air brakes were the predominant form of train braking in most countries during the steam era.
The primary competitor to the air brake was the vacuum brake, in which a steam-operated ejector is mounted on the engine instead of the air pump, to create the vacuum required to power the brake system. A secondary ejector or crosshead vacuum pump is used to maintain the vacuum in the system against the small leaks in the pipe connections between carriages and wagons. Vacuum brakes were the predominant form of train braking in the United Kingdom and countries that adopted its practices, such as India and South Africa, during the steam era.
Steam locomotives are fitted with sandboxes from which sand can be deposited on top of the rail to improve traction and braking in wet or icy weather. On American locomotives, the sandboxes, or sand domes, are usually mounted on top of the boiler. In Britain, the limited loading gauge precludes this, so the sandboxes are mounted just above, or just below, the running plate.
Lubrication
The pistons and valves on the earliest locomotives were lubricated by the enginemen dropping a lump of tallow down the blast pipe. More sophisticated methods of delivering the substance were soon developed. Tallow adheres well to cylinder walls and is more effective than mineral oil in resisting the action of water. It remains a constituent of modern steam cylinder oil formulation.
As speeds and distances increased, mechanisms were developed that injected thick mineral oil into the steam supply. The first, a displacement lubricator, mounted in the cab, uses a controlled stream of steam condensing into a sealed container of oil. Water from the condensed steam displaces the oil into pipes. The apparatus is usually fitted with sight-glasses to confirm the rate of supply. A later method uses a mechanical pump worked from one of the crossheads. In both cases, the supply of oil is proportional to the speed of the locomotive.
Lubricating the frame components (axle bearings, horn blocks and bogie pivots) depends on capillary action: trimmings of worsted yarn are trailed from oil reservoirs into pipes leading to the respective component. The rate of oil supplied is controlled by the size of the bundle of yarn and not the speed of the locomotive, so it is necessary to remove the trimmings (which are mounted on wire) when stationary. However, at regular stops (such as a terminating station platform), oil finding its way onto the track can still be a problem.
Crankpin and crosshead bearings carry small cup-shaped reservoirs for oil. These have feed pipes to the bearing surface that start above the normal fill level, or are kept closed by a loose-fitting pin, so that only when the locomotive is in motion does oil enter. In United Kingdom practice, the cups are closed with simple corks, but these have a piece of porous cane pushed through them to admit air. It is customary for a small capsule of pungent oil (aniseed or garlic) to be incorporated in the bearing metal to warn if the lubrication fails and excess heating or wear occurs.
Blower
When the locomotive is running under power, a draught on the fire is created by the exhaust steam directed up the chimney by the blastpipe. Without draught, the fire will quickly die down and steam pressure will fall. When the locomotive is stopped, or coasting with the regulator closed, there is no exhaust steam to create a draught, so the draught is maintained by means of a blower. This is a ring placed either around the base of the chimney, or around the blast pipe orifice, containing several small steam nozzles directed up the chimney. These nozzles are fed with steam directly from the boiler, controlled by the blower valve. When the regulator is open, the blower valve is closed; when the driver intends to close the regulator, he will first open the blower valve. It is important that the blower be opened before the regulator is closed, since without draught on the fire, there may be backdraught – where atmospheric air blows down the chimney, causing the flow of hot gases through the boiler tubes to be reversed, with the fire itself being blown through the firehole onto the footplate, with serious consequences for the crew. The risk of backdraught is higher when the locomotive enters a tunnel because of the pressure shock. The blower is also used to create draught when steam is being raised at the start of the locomotive's duty, at any time when the driver needs to increase the draught on the fire, and to clear smoke from the driver's line of vision.
Blowbacks were fairly common. In a 1955 report on an accident near Dunstable, the Inspector wrote, "In 1953 twenty-three cases, which were not caused by an engine defect, were reported and they resulted in 26 enginemen receiving injuries. In 1954, the number of occurrences and of injuries were the same and there was also one fatal casualty." They remain a problem, as evidenced by the 2012 incident with BR Standard Class 7 70013 Oliver Cromwell.
Buffers
In British and European (except former Soviet Union countries) practice, locomotives usually have buffers at each end to absorb compressive loads ("buffets"). The tensional load of drawing the train (draft force) is carried by the coupling system. Together these control slack between the locomotive and train, absorb minor impacts and provide a bearing point for pushing movements.
In Canadian and American practice, all of the forces between the locomotive and cars are handled through the coupler – particularly the Janney coupler, long standard on American railroad rolling stock – and its associated draft gear, which allows some limited slack movement. Small dimples called "poling pockets" at the front and rear corners of the locomotive allowed cars to be pushed onto an adjacent track using a pole braced between the locomotive and the cars. In Britain and Europe, North American style "buckeye" and other couplers that handle forces between items of rolling stock have become increasingly popular.
Pilots
A pilot was usually fixed to the front end of locomotives, although in European and a few other railway systems including New South Wales, they were considered unnecessary. Plough-shaped, sometimes called "cow catchers", they were quite large and were designed to remove obstacles from the track such as cattle, bison, other animals or tree limbs. Though unable to "catch" stray cattle, these distinctive items remained on locomotives until the end of steam. Switching engines usually replaced the pilot with small steps, known as footboards. Many systems used the pilot and other design features to produce a distinctive appearance.
Headlights
When night operations began, railway companies in some countries equipped their locomotives with lights to allow the driver to see what lay ahead of the train, or to enable others to see the locomotive. Headlights were originally oil or acetylene lamps, but when electric arc lamps became available in the late 1880s, they quickly replaced the older types.
Britain did not adopt bright headlights as they would affect night vision and so could mask the low-intensity oil lamps used in the semaphore signals and at each end of trains, increasing the danger of missing signals, especially on busy tracks. Locomotive stopping distances were also normally much greater than the range of headlights, and the railways were well-signalled and fully fenced to prevent livestock and people from straying onto them, largely negating the need for bright lamps. Thus low-intensity oil lamps continued to be used, positioned on the front of locomotives to indicate the class of each train. Four "lamp irons" (brackets on which to place the lamps) were provided: one below the chimney and three evenly spaced across the top of the buffer beam. The exception to this was the Southern Railway and its constituents, who added an extra lamp iron each side of the smokebox, and the arrangement of lamps (or in daylight, white circular plates) told railway staff the origin and destination of the train. On all vehicles, equivalent lamp irons were also provided on the rear of the locomotive or tender for when the locomotive was running tender- or bunker-first.
In some countries, heritage steam operation continues on the national network. Some railway authorities have mandated powerful headlights on at all times, including during daylight. This was to further inform the public or track workers of any active trains.
Bells and whistles
Locomotives used bells and steam whistles from earliest days of steam locomotion. In the United States, India and Canada, bells warned of a train in motion. In Britain, where all lines are by law fenced throughout, bells were only a requirement on railways running on a road (i.e. not fenced off), for example a tramway along the side of the road or in a dockyard. Consequently, only a minority of locomotives in the UK carried bells. Whistles are used to signal personnel and give warnings. Depending on the terrain the locomotive was being used in, the whistle could be designed for long-distance warning of impending arrival, or for more localised use.
Early bells and whistles were sounded through pull-string cords and levers. Automatic bell ringers came into widespread use in the US after 1910.
Automatic control
From the early 20th century operating companies in such countries as Germany and Britain began to fit locomotives with Automatic Warning System (AWS) in-cab signalling, which automatically applied the brakes when a signal was passed at "caution". In Britain, these became mandatory in 1956. In the United States, the Pennsylvania Railroad also fitted their locomotives with such devices.
Booster engines
The booster engine was an auxiliary steam engine which provided extra tractive effort for starting. It was a low-speed device, usually mounted on the trailing truck. It was disengaged via an idler gear at a low speed, e.g. 30 km/h. Boosters were widely used in the US and tried experimentally in Britain and France. On the narrow-gauged New Zealand railway system, six Kb 4-8-4 locomotives were fitted with boosters, the only gauge engines in the world to have such equipment.
Booster engines were also fitted to tender trucks in the US and known as auxiliary locomotives. Two and even three truck axles were connected together using side rods which limited them to slow-speed service.
Firedoor
The firedoor is used to cover the firehole when coal is not being added. It serves two purposes, first, it prevents air being drawn over the top of the fire, rather forcing it to be drawn through it. The second purpose is to safeguard the train crew against blowbacks. It does, however, have a means to allow some air to pass over the top of the fire (referred to as "secondary air") to complete the combustion of gases produced by the fire.
Firedoors come in multiple designs, the most basic of which is a single piece which is hinged on one side and can swing open onto the footplate. This design has two issues. First, it takes up much room on the footplate, and second, the draught will tend to pull it completely shut, thus cutting off any secondary air. To compensate for this some locomotives are fitted with a latch that prevents the firedoor from closing completely whereas others have a small vent on the door that may be opened to allow secondary air to flow through. Though it was considered to design a firedoor that opens inwards into the firebox thus preventing the inconvenience caused on the footplate, such a door would be exposed to the full heat of the fire and would likely deform, thus becoming useless.
A more popular type of firedoor consists of a two-piece sliding door operated by a single lever. There are tracks above and below the firedoor which the door runs along. These tracks are prone to becoming jammed by debris and the doors required more effort to open than the aforementioned swinging door. In order to address this some firedoors use powered operation which utilized a steam or air cylinder to open the door. Among these are the butterfly doors which pivot at the upper corner, the pivoting action offers low resistance to the cylinder that opens the door.
Variations
Numerous variations on the basic locomotive occurred as railways attempted to improve efficiency and performance.
Cylinders
Early steam locomotives had two cylinders, one either side, and this practice persisted as the simplest arrangement. The cylinders could be mounted between the mainframes (known as "inside" cylinders), or mounted outside the frames and driving wheels ("outside" cylinders). Inside cylinders drive cranks built into the driving axle; outside cylinders drive cranks on extensions to the driving axles.
Later designs employed three or four cylinders, mounted both inside and outside the frames, for a more even power cycle and greater power output. This was at the expense of more complicated valve gear and increased maintenance requirements. In some cases the third cylinder was added inside simply to allow for smaller diameter outside cylinders, and hence reduce the width of the locomotive for use on lines with a restricted loading gauge, for example the SR K1 and U1 classes.
Most British express-passenger locomotives built between 1930 and 1950 were or types with three or four cylinders (e.g. GWR 6000 Class, LMS Coronation Class, SR Merchant Navy Class, LNER Gresley Class A3). From 1951, all but one of the 999 new British Rail standard class steam locomotives across all types used 2-cylinder configurations for easier maintenance.
Valve gear
Early locomotives used a simple valve gear that gave full power in either forward or reverse. Soon the Stephenson valve gear allowed the driver to control cut-off; this was largely superseded by Walschaerts valve gear and similar patterns. Early locomotive designs using slide valves and outside admission were relatively easy to construct, but inefficient and prone to wear. Eventually, slide valves were superseded by inside admission piston valves, though there were attempts to apply poppet valves (commonly used in stationary engines) in the 20th century. Stephenson valve gear was generally placed within the frame and was difficult to access for maintenance; later patterns applied outside the frame were more readily visible and maintained.
Compounding
Compound locomotives were used from 1876, expanding the steam twice or more through separate cylinders – reducing thermal losses caused by cylinder cooling. Compound locomotives were especially useful in trains where long periods of continuous efforts were needed. Compounding contributed to the dramatic increase in power achieved by André Chapelon's rebuilds from 1929. A common application was in articulated locomotives, the most common being that designed by Anatole Mallet, in which the high-pressure stage was attached directly to the boiler frame; in front of this was pivoted a low-pressure engine on its own frame, which takes the exhaust from the rear engine.
Articulated locomotives
Very powerful locomotives tend to be longer than those with lower power output, but long rigid-framed designs are impracticable for the tight curves frequently found on narrow-gauge railways. Various designs for articulated locomotives were developed to overcome this problem. The Mallet and the Garratt were the two most popular. They had a single boiler and two engine units (sets of cylinders and driving wheels): both of the Garratt's engine units were on swivelling frames, whereas one of the Mallet's was on a swivelling frame and the other was fixed under the boiler unit. A few triplex locomotives were also designed, with a third engine unit under the tender. Other less common variations included the Fairlie locomotive, which had two boilers back-to-back on a common frame, with two separate engine units.
Duplex types
Duplex locomotives, containing two engines in one rigid frame, were also tried, but were not notably successful. For example, the Pennsylvania Railroad class T1, designed for very fast running, suffered recurring and ultimately unfixable slippage problems throughout their careers.
Geared locomotives
For locomotives where a high starting torque and low speed were required, the conventional direct drive approach was inadequate. "Geared" steam locomotives, such as the Shay, the Climax and the Heisler, were developed to meet this need on industrial, logging, mine and quarry railways. The common feature of these three types was the provision of reduction gearing and a drive shaft between the crankshaft and the driving axles. This arrangement allowed the engine to run at a much higher speed than the driving wheels compared to the conventional design, where the ratio is 1:1.
Cab forward
In the United States on the Southern Pacific Railroad, a series of cab forward locomotives were produced with the cab and the firebox at the front of the locomotive and the tender behind the smokebox, so that the engine appeared to run backwards. This was only possible by using oil-firing. Southern Pacific selected this design to provide air free of smoke for the engine driver to breathe as the locomotive passed through mountain tunnels and snow sheds. Another variation was the Camelback locomotive, with the cab situated halfway along the boiler. In England, Oliver Bulleid developed the SR Leader class locomotive during the nationalisation process in the late 1940s. The locomotive was heavily tested but several design faults (such as coal firing and sleeve valves) meant that this locomotive and the other part-built locomotives were scrapped. The cab-forward design was taken by Bulleid to Ireland, where he moved after nationalisation, where he developed the "turfburner". This locomotive was more successful, but was scrapped due to the dieselisation of the Irish railways.
The only preserved cab forward locomotive is Southern Pacific 4294 in Sacramento, California.
In France, the three Heilmann locomotives were built with a cab forward design.
Steam turbines
Steam turbines were created as an attempt to improve the operation and efficiency of steam locomotives. Experiments with steam turbines using direct-drive and electrical transmissions in various countries proved mostly unsuccessful. The London, Midland & Scottish Railway built the Turbomotive, a largely successful attempt to prove the efficiency of steam turbines. Had it not been for the outbreak of World War II, more may have been built. The Turbomotive ran from 1935 to 1949, when it was rebuilt into a conventional locomotive because many parts required replacement, an uneconomical proposition for a "one-off" locomotive. In the United States, Union Pacific, Chesapeake & Ohio and Norfolk & Western (N&W) railways all built turbine-electric locomotives. The Pennsylvania Railroad (PRR) also built turbine locomotives, but with a direct-drive gearbox. However, all designs failed due to dust, vibration, design flaws or inefficiency at lower speeds. The final one remaining in service was the N&W's, retired in January 1958. The only truly successful design was the TGOJ MT3, used for hauling iron ore from Grängesberg in Sweden to the ports of Oxelösund. Despite functioning correctly, only three were built. Two of them are preserved in working order in museums in Sweden.
Fireless locomotive
In a fireless locomotive the boiler is replaced by a steam accumulator, which is charged with steam (actually water at a temperature well above boiling point, () from a stationary boiler. Fireless locomotives were used where there was a high fire risk (e.g. oil refineries), where cleanliness was important (e.g. food-production plants) or where steam is readily available (e.g. paper mills and power stations where steam is either a by-product or is cheaply available). The water vessel ("boiler") is heavily insulated, the same as with a fired locomotive. Until all the water has boiled away, the steam pressure does not drop except as the temperature drops.
Another class of fireless locomotive is a compressed-air locomotive.
Mixed power
Steam diesel hybrid locomotive
Mixed power locomotives, utilising both steam and diesel propulsion, have been produced in Russia, Britain and Italy.
Electric–steam locomotive
Under unusual conditions (lack of coal, abundant hydroelectricity) some locomotives in Switzerland were modified to use electricity to heat the boiler, making them electric–steam locomotives.
Steam-electric locomotive
A steam-electric locomotive uses electric transmission, like diesel-electric locomotives, except that a steam engine instead of a diesel engine is used to drive a generator. Three such locomotives were built by the French engineer in the 1890s.
Categorisation
Steam locomotives are categorised by their wheel arrangement. The two dominant systems for this are the Whyte notation and UIC classification.
The Whyte notation, used in most English-speaking and Commonwealth countries, represents each set of wheels with a number. These numbers typically represented the number of unpowered leading wheels, followed by the number of driving wheels (sometimes in several groups), followed by the number of un-powered trailing wheels. For example, a yard engine with only 4 driven wheels, without any leading or trailing wheels, would be categorised as a wheel arrangement. A locomotive with a 4-wheel leading truck, followed by 6 drive wheels, and a 2-wheel trailing truck, would be classed as a . Different arrangements were given names which usually reflect the first usage of the arrangement; for instance, the "Santa Fe" type () is so called because the first examples were built for the Atchison, Topeka and Santa Fe Railway. These names were informally given and varied according to region and even politics.
The UIC classification is used mostly in European countries apart from the United Kingdom. It designates consecutive pairs of wheels (informally "axles") with a number for non-driving wheels and a capital letter for driving wheels (A=1, B=2, etc.) So a Whyte designation would be an equivalent to a 2-C-1 UIC designation.
On many railroads, locomotives were organised into classes. These broadly represented locomotives which could be substituted for each other in service, but most commonly a class represented a single design. As a rule classes were assigned some sort of code, generally based on the wheel arrangement. Classes also commonly acquired nicknames, such as Pug (a small shunting locomotive), representing notable (and sometimes uncomplimentary) features of the locomotives.
Performance
Measurement
In the steam locomotive era, two measures of locomotive performance were generally applied. At first, locomotives were rated by tractive effort, defined as the average force developed during one revolution of the driving wheels at the railhead. This can be roughly calculated by multiplying the total piston area by 85% of the boiler pressure (a rule of thumb reflecting the slightly lower pressure in the steam chest above the cylinder) and dividing by the ratio of the driver diameter over the piston stroke. However, the precise formula is
where is the bore of the cylinder (diameter) in inches,
is the cylinder stroke, in inches,
is boiler pressure in pounds per square inch,
is the diameter of the driving wheel in inches,
and is a factor that depends on the effective cut-off. In the US, is usually set at 0.85, but lower on engines that have maximum cutoff limited to 50–75%.
The tractive effort is only the "average" force, as not all effort is constant during the one revolution of the drivers. At some points of the cycle, only one piston is exerting turning moment and at other points, both pistons are working. Not all boilers deliver full power at starting, and the tractive effort also decreases as the rotating speed increases.
Tractive effort is a measure of the heaviest load a locomotive can start or haul at very low speed over the ruling grade in a given territory. However, as the pressure grew to run faster goods and heavier passenger trains, tractive effort was seen to be an inadequate measure of performance because it did not take into account speed. Therefore, in the 20th century, locomotives began to be rated by power output. A variety of calculations and formulas were applied, but in general railways used dynamometer cars to measure tractive force at speed in actual road testing.
British railway companies have been reluctant to disclose figures for drawbar horsepower and have usually relied on continuous tractive effort instead.
Relation to wheel arrangement
Classification is indirectly connected to locomotive performance. Given adequate proportions of the rest of the locomotive, power output is determined by the size of the fire, and for a bituminous coal-fuelled locomotive, this is determined by the grate area. Modern non-compound locomotives are typically able to produce about 40 drawbar horsepower per square foot of grate. Tractive force, as noted earlier, is largely determined by the boiler pressure, the cylinder proportions and the size of the driving wheels. However, it is also limited by the weight on the driving wheels (termed "adhesive weight"), which needs to be at least four times the tractive effort.
The weight of the locomotive is roughly proportional to the power output; the number of axles required is determined by this weight divided by the axleload limit for the trackage where the locomotive is to be used. The number of driving wheels is derived from the adhesive weight in the same manner, leaving the remaining axles to be accounted for by the leading and trailing bogies. Passenger locomotives conventionally had two-axle leading bogies for better guidance at speed; on the other hand, the vast increase in the size of the grate and firebox in the 20th century meant that a trailing bogie was called upon to provide support. In Europe, some use was made of several variants of the Bissel bogie in which the swivelling movement of a single axle truck controls the lateral displacement of the front driving axle (and in one case the second axle too). This was mostly applied to 8-coupled express and mixed traffic locomotives, and considerably improved their ability to negotiate curves whilst restricting overall locomotive wheelbase and maximising adhesion weight.
As a rule, shunting engines (US: switching engines) omitted leading and trailing bogies, both to maximise tractive effort available and to reduce wheelbase. Speed was unimportant; making the smallest engine (and therefore smallest fuel consumption) for the tractive effort was paramount. Driving wheels were small and usually supported the firebox as well as the main section of the boiler. Banking engines (US: helper engines) tended to follow the principles of shunting engines, except that the wheelbase limitation did not apply, so banking engines tended to have more driving wheels. In the US, this process eventually resulted in the Mallet type engine with its many driven wheels, and these tended to acquire leading and then trailing bogies as guidance of the engine became more of an issue.
As locomotive types began to diverge in the late 19th century, freight engine designs at first emphasised tractive effort, whereas those for passenger engines emphasised speed. Over time, freight locomotive size increased, and the overall number of axles increased accordingly; the leading bogie was usually a single axle, but a trailing truck was added to larger locomotives to support a larger firebox that could no longer fit between or above the driving wheels. Passenger locomotives had leading bogies with two axles, fewer driving axles, and very large driving wheels in order to limit the speed at which the reciprocating parts had to move.
In the 1920s, the focus in the United States turned to horsepower, epitomised by the "super power" concept promoted by the Lima Locomotive Works, although tractive effort was still the prime consideration after World War I to the end of steam. Goods trains were designed to run faster, while passenger locomotives needed to pull heavier loads at speed. This was achieved by increasing the size of grate and firebox without changes to the rest of the locomotive, requiring the addition of a second axle to the trailing truck. Freight s became s while s became s. Similarly, passenger s became s. In the United States this led to a convergence on the dual-purpose and the articulated configuration, which was used for both freight and passenger service. Mallet locomotives went through a similar transformation, evolving from bank engines into huge mainline locomotives with much larger fireboxes; their driving wheels were also increased in size in order to allow faster running.
Manufacture
Most-manufactured classes
The most-manufactured single class of steam locomotive in the world is the Russian locomotive class E steam locomotive with around 11,000 produced both in Russia and other countries such as Czechoslovakia, Germany, Sweden, Hungary and Poland. The Russian locomotive class O numbered 9,129 locomotives, built between 1890 and 1928. Around 7,000 units were produced of the German DRB Class 52 Kriegslok.
In Britain, 863 of the GWR 5700 Class were built, and 943 of the DX class of the London and North Western Railway including 86 engines built for the Lancashire and Yorkshire Railway.
United Kingdom
Before the 1923 Grouping Act, production in the UK was mixed. The larger railway companies built locomotives in their own workshops, with the smaller ones and industrial concerns ordering them from outside builders. A large market for outside builders existed due to the home-build policy exercised by the main railway companies. An example of a pre-grouping works was the one at Melton Constable, which maintained and built some of the locomotives for the Midland and Great Northern Joint Railway. Other works included one at Boston (an early GNR building) and Horwich Works.
Between 1923 and 1947, the Big Four railway companies (the Great Western Railway, the London, Midland & Scottish Railway, the London & North Eastern Railway and the Southern Railway) all built most of their own locomotives, only buying locomotives from outside builders when their own works were fully occupied (or as a result of government-mandated standardisation during wartime).
From 1948, British Railways (BR) allowed the former Big Four companies (now designated as "Regions") to continue to produce their own designs, but also created a range of standard locomotives which supposedly combined the best features from each region. Although a policy of dieselisation was adopted in 1955, BR continued to build new steam locomotives until 1960, with the final engine being named Evening Star.
Some independent manufacturers produced steam locomotives for a few more years, with the last British-built industrial steam locomotive being constructed by Hunslet in 1971. Since then, a few specialised manufacturers have continued to produce small locomotives for narrow gauge and miniature railways, but as the prime market for these is the tourist and heritage railway sector, the demand for such locomotives is limited. In November 2008, a new build main line steam locomotive, 60163 Tornado, was tested on UK mainlines for eventual charter and tour use.
Sweden
In the 19th and early 20th centuries, most Swedish steam locomotives were manufactured in Britain. Later, however, most steam locomotives were built by local factories including NOHAB in Trollhättan and ASJ in Falun. One of the most successful types was the class "B" (), inspired by the Prussian class P8. Many of the Swedish steam locomotives were preserved during the Cold War in case of war. During the 1990s, these steam locomotives were sold to non-profit associations or abroad, which is why the Swedish class B, class S () and class E2 () locomotives can now be seen in Britain, the Netherlands, Germany and Canada.
United States
Locomotives for American railroads were nearly always built in the United States with very few imports, except in the earliest days of steam engines. This was due to the basic differences of markets in the United States which initially had many small markets located large distances apart, in contrast to Europe's higher density of markets. Locomotives that were cheap and rugged and could go large distances over cheaply built and maintained tracks were required. Once the manufacture of engines was established on a wide scale there was very little advantage to buying an engine from overseas that would have to be customised to fit the local requirements and track conditions. Improvements in engine design of both European and US origin were incorporated by manufacturers when they could be justified in a generally very conservative and slow-changing market. With the notable exception of the USRA standard locomotives built during World War I, in the United States, steam locomotive manufacture was always semi-customised. Railroads ordered locomotives tailored to their specific requirements, though some basic design features were always present. Railroads developed some specific characteristics; for example, the Pennsylvania Railroad and the Great Northern Railway had a preference for the Belpaire firebox. In the United States, large-scale manufacturers constructed locomotives for nearly all rail companies, although nearly all major railroads had shops capable of heavy repairs and some railroads (for example, the Norfolk and Western Railway and the Pennsylvania Railroad, which had two erecting shops) constructed locomotives entirely in their own shops. Companies manufacturing locomotives in the US included Baldwin Locomotive Works, American Locomotive Company (ALCO), and Lima Locomotive Works. Altogether, between 1830 and 1950, over 160,000 steam locomotives were built in the United States, with Baldwin accounting for the largest share, nearly 70,000.
Steam locomotives required regular and, compared to a diesel-electric engine, frequent service and overhaul (often at government-regulated intervals in Europe and the US). Alterations and upgrades regularly occurred during overhauls. New appliances were added, unsatisfactory features removed, cylinders improved or replaced. Almost any part of the locomotive, including boilers, was replaced or upgraded. When service or upgrades got too expensive the locomotive was traded off or retired. On the Baltimore and Ohio Railroad two locomotives were dismantled; the boilers were placed onto two new Class T locomotives and the residual wheel machinery made into a pair of Class U switchers with new boilers. Union Pacific's fleet of 3-cylinder engines were converted into two-cylinder engines in 1942, because of high maintenance problems.
Australia
In Sydney, Clyde Engineering and the Eveleigh Railway Workshops both built steam locomotives for the New South Wales Government Railways. These include the C38 class ; the first five were built at Clyde with streamlining, the other 25 locomotives were built at Eveleigh (13) and Cardiff Workshops (12) near Newcastle. In Queensland, steam locomotives were locally constructed by Walkers. Similarly, the South Australian Railways also manufactured steam locomotives locally at Islington Railway Workshops in Adelaide. Victorian Railways constructed most of their locomotives at its Newport Workshops and in Bendigo, while in the early days locomotives were built at the Phoenix Foundry in Ballarat. Locomotives constructed at the Newport shops ranged from the nA class built for the narrow gauge, up to the H class – the largest conventional locomotive ever to operate in Australia, weighing 260 tons. However, the title of largest locomotive ever used in Australia goes to the 263-ton New South Wales AD60 class locomotive Garratt, built by Beyer, Peacock & Company in England. Most steam locomotives used in Western Australia were built in the United Kingdom, though some examples were designed and built locally at the Western Australian Government Railways' Midland Railway Workshops. The 10 WAGR S class locomotives (introduced in 1943) were the only class of steam locomotive to be wholly conceived, designed and built in Western Australia, while the Midland workshops notably participated in the Australia-wide construction program of Australian Standard Garratts – these wartime locomotives were built at Midland in Western Australia, Clyde Engineering in New South Wales, Newport in Victoria and Islington in South Australia and saw varying degrees of service in all Australian states.
The end of steam in general use
The introduction of electric locomotives around the turn of the 20th century and later diesel-electric locomotives spelled the beginning of a decline in the use of steam locomotives, although it was some time before they were phased out of general use. As diesel power (especially with electric transmission) became more reliable in the 1930s, it gained a foothold in North America. The full transition away from steam power in North America took place during the 1950s. In continental Europe, large-scale electrification had replaced steam power by the 1970s. Steam was a familiar technology, adapted well to local facilities, and also consumed a wide variety of fuels; this led to its continued use in many countries until the end of the 20th century.
Steam engines have considerably less thermal efficiency than modern diesels, requiring constant maintenance and labour to keep them operational. Water is required at many points throughout a rail network, making it a major problem in desert areas, as are found in some regions of the United States, Australia and South Africa. In places where water is available, it may be hard, which can cause "scale" to form, composed mainly of calcium carbonate, magnesium hydroxide and calcium sulfate. Calcium and magnesium carbonates tend to be deposited as off-white solids on the inside the surfaces of pipes and heat exchangers. This precipitation is principally caused by thermal decomposition of bicarbonate ions but also happens in cases where the carbonate ion is at saturation concentration. The resulting build-up of scale restricts the flow of water in pipes. In boilers, the deposits impair the flow of heat into the water, reducing the heating efficiency and allowing the metal boiler components to overheat.
The reciprocating mechanism on the driving wheels of a two-cylinder single expansion steam locomotive tended to pound the rails (see hammer blow), thus requiring more maintenance. Raising steam from coal took a matter of hours, and created serious pollution problems. Coal-burning locomotives required fire cleaning and ash removal between turns of duty. Diesel or electric locomotives, by comparison, drew benefit from new custom-built servicing facilities. The smoke from steam locomotives was also deemed objectionable; the first electric and diesel locomotives were developed in response to smoke abatement requirements, although this did not take into account the high level of less-visible pollution in diesel exhaust smoke, especially when idling. In some countries, however, power for electric locomotives is derived from steam generated in power stations, which are often run by coal.
Revival
Dramatic increases in the cost of diesel fuel prompted several initiatives to revive steam power. However, none of these has progressed to the point of production and, as of the early 21st century, steam locomotives operate only in a few isolated regions of the world and in tourist operations.
As early as 1975, railway enthusiasts in the United Kingdom began building new steam locomotives. That year, Trevor Barber completed his gauge locomotive Trixie which ran on the Meirion Mill Railway. From the 1990s onwards, the number of new builds being completed rose dramatically with new locos completed by the narrow-gauge Ffestiniog and Corris railways in Wales. The Hunslet Engine Company was revived in 2005, and began building steam locomotives on a commercial basis. A standard-gauge LNER Peppercorn Pacific "Tornado" was completed at Hopetown Works, Darlington, and made its first run on 1 August 2008. It entered main line service later in 2008. over half-a-dozen projects to build working replicas of extinct steam engines are going ahead, in many cases using existing parts from other types to build them. Examples include BR 72010 Hengist, BR Class 3MT No. 82045, BR Class 2MT No. 84030, Brighton Atlantic Beachy Head, the LMS 5551 The Unknown Warrior project, GWR "47xx 4709, 2999 Lady of Legend, 1014 County of Glamorgan and 6880 Betton Grange projects. These United Kingdom based new build projects are further complemented by the new build Pennsylvania Railroad 5550 project in the United States. One of the group's goals is to surpass the steam locomotive speed record held by the 4468 Mallard when the 5550 is completed and for the 5550 to fill in a huge gap in steam locomotive preservation.
In 1980, American financier Ross Rowland established American Coal Enterprises to develop a modernised coal-fired steam locomotive. His ACE 3000 concept attracted considerable attention, but was never built.
In 1998, in his book The Red Devil and Other Tales from the Age of Steam, David Wardale put forward the concept of a high-speed high-efficiency "Super Class 5 4-6-0" locomotive for future steam haulage of tour trains on British main lines. The idea was formalised in 2001 by the formation of 5AT Project dedicated to developing and building the 5AT Advanced Technology Steam Locomotive, but it never received any major railway backing.
Locations where new builds are taking place include:
GWR 1014 County of Glamorgan & GWR 2999 Lady of Legend, both being built at Didcot Railway Centre
GWR 6880 Betton Grange, GWR 4709 & LMS 5551 The Unknown Warrior, all being built at Llangollen Railway
LNER 2007 Prince of Wales, Darlington Locomotive Works
LNER 2001 Cock O' The North, Doncaster
Pennsylvania Railroad 5550, Pottstown, Pennsylvania
BR 72010 Hengist, Great Central Railway
BR 77021, TBA
BR 82045, Severn Valley Railway
BR 84030 & LBSCR 32424 Beachy Head, both being built at Bluebell Railway
MS&LR/GCR 567, Ruddington Great Central Railway, Northern Section
VR V499, Victoria, Australia
In 2012, the Coalition for Sustainable Rail project was started in the US with the goal of creating a modern higher-speed steam locomotive, incorporating the improvements proposed by Livio Dante Porta and others, and using torrefied biomass as solid fuel. The fuel has been recently developed by the University of Minnesota in a collaboration between the university's Institute on the Environment (IonE) and Sustainable Rail International (SRI), an organisation set up to explore the use of steam traction in a modern railway setup. The group have received the last surviving (but non-running) ATSF 3460 class steam locomotive (No. 3463) via donation from its previous owner in Kansas, the Great Overland Station Museum. They hope to use it as a platform for developing "the world's cleanest, most powerful passenger locomotive", capable of speeds up to . Named "Project 130", it aims to break the world steam-train speed record set by LNER Class A4 4468 Mallard in the UK at . However, any demonstration of the project's claims is yet to be seen.
In Germany, a small number of fireless steam locomotives are still working in industrial service, e.g. at power stations, where an on-site supply of steam is readily available.
The small town of Wolsztyn, Poland, approximately from the historic city of Poznań, is the last place in the world where one can ride a regularly scheduled passenger train pulled by steam power. The locomotive shed at Wolsztyn is the last of its kind in the world. There are several working locomotives that haul daily commuter service between Wolsztyn, Poznan, Leszo and other neighboring cities. One can partake in footplate courses via The Wolsztyn Experience. There is no place left in the world that still operates daily, non-tourist steam powered commuter/passenger service other than here at Wolsztyn. There are several Polish-built OL49-class 2-6-2 general purpose locomotives and one PT47 class 2-8-2 in regular service. Each May, Wolsztyn is the site of a steam locomotive festival which brings visiting locomotives - often well over a dozen each year all operating. These operations are not done for tourism or museum/historical purposes; this is the last non-diesel rail line on the PKP (Polish State Network) that has been converted to diesel power.
The Swiss company Dampflokomotiv- und Maschinenfabrik DLM AG delivered eight steam locomotives to rack railways in Switzerland and Austria between 1992 and 1996. Four of them are now the main traction on the Brienz Rothorn Bahn; the four others were built for the Schafbergbahn in Austria, where they run 90% of the trains.
The same company also rebuilt a German DR Class 52.80 2-10-0 locomotive to new standards with modifications such as roller bearings, light oil firing and boiler insulation.
Climate change
The future use of steam locomotives in the United Kingdom is in doubt because of government policy on climate change. The Heritage Railway Association is working with the All-Party Parliamentary Group on Heritage Rail in an effort to continue running steam locomotives on coal.
Many tourist railroads use oil-fired steam locomotives (or have converted their locomotives to run on oil) to reduce their environmental footprint, and because fuel oil can be easier to obtain than coal of the proper type and sizing for locomotives. For example, the Grand Canyon Railway runs its steam locomotives on used vegetable oil.
An organization called the Coalition for Sustainable Rail (CSR) is developing an environmentally friendly coal substitute made from torrefied biomass. In early 2019, they performed a series of tests using Everett Railroad to evaluate the performance of the biofuel, with positive results. The biofuel was found to burn slightly faster and hotter than coal. The goal of the project is primarily to find a sustainable fuel for historic steam locomotives on tourist railroads, but CSR has also suggested that, in the future, steam locomotives powered by torrefied biomass could be an environmentally and economically superior alternative to diesel locomotives. Also, a large vat containing salt may be used without needing to replenishing the medium. Large heating elements would be one method of recharging the system, however, it is possible to pump molten salt as well, removing the cooled salt and replenishing from facilities which contain a much larger vat.
Steam locomotives in popular culture
Steam locomotives have been present in popular culture since the 19th century. Folk songs from that period including "I've Been Working on the Railroad" and the "Ballad of John Henry" are a mainstay of American music and culture.
Many steam locomotive toys have been made, and railway modelling is a popular hobby.
Steam locomotives are often portrayed in fictional works, notably The Railway Series by the Rev W. V. Awdry, The Little Engine That Could by Watty Piper, The Polar Express by Chris Van Allsburg, and the Hogwarts Express from J.K. Rowling's Harry Potter series. They have also been featured in many children's television shows, such as Thomas & Friends, based on characters from the books by Awdry, and Ivor the Engine created by Oliver Postgate.
The Hogwarts Express also appears in the Harry Potter series of films, portrayed by GWR 4900 Class 5972 Olton Hall in a special Hogwarts livery. The Polar Express appears in the animated movie of the same name.
An elaborate, themed funicular Hogwarts Express ride is featured in the Universal Orlando Resort in Florida, connecting the Harry Potter section of Universal Studios with the Islands of Adventure theme park.
The Polar Express is recreated on many heritage railroads in the United States, including the North Pole Express pulled by the Pere Marquette 1225 locomotive, which is operated by the Steam Railroading Institute in Owosso, Michigan. According to author Van Allsburg, this locomotive was the inspiration for the story and it was used in the production of the movie.
A number of computer and video games feature steam locomotives, such as Railroad Tycoon and Railway Empire. Railroad Tycoon, produced in 1990, was named "one of the best computer games of the year".
There are two notable examples of steam locomotives used as charges on heraldic coats of arms. One is that of Darlington, which displays Locomotion No. 1. The other is the original coat of arms of Swindon, not currently in use, which displays a basic steam locomotive.
Steam locomotives are a popular topic for coin collectors. The 1950 Silver 5 Peso coin of Mexico has a steam locomotive on its reverse as the prominent feature.
The 20 euro Biedermeier Period coin, minted 11 June 2003, shows on the obverse an early model steam locomotive (the Ajax) on Austria's first railway line, the Kaiser Ferdinands-Nordbahn. The Ajax can still be seen today in the Technisches Museum Wien.
As part of the 50 State Quarters program, the quarter representing the US state of Utah depicts the ceremony where the two halves of the First transcontinental railroad met at Promontory Summit in 1869. The coin recreates a popular image from the ceremony with steam locomotives from each company facing each other while the golden spike is being driven.
The novel "Night on the Galactic Railroad" by Kenji Miyazawa is centered on the idea of a steam train traveling among the stars. Miyazawa's novel later inspired Leiji Matsumoto's successful "Galaxy Express 999" series.
Another Japanese televisual franchise, Super Sentai, features monsters based on steam locomotives.
Charge Man, a Robot Master from the fifth installment of the Mega Man series is based on a steam locomotive.
| Technology | Trains | null |
196974 | https://en.wikipedia.org/wiki/Cough | Cough | A cough is a sudden expulsion of air through the large breathing passages which can help clear them of fluids, irritants, foreign particles and microbes. As a protective reflex, coughing can be repetitive with the cough reflex following three phases: an inhalation, a forced exhalation against a closed glottis, and a violent release of air from the lungs following opening of the glottis, usually accompanied by a distinctive sound.
Frequent coughing usually indicates the presence of a disease. Many viruses and bacteria benefit, from an evolutionary perspective, by causing the host to cough, which helps to spread the disease to new hosts. Irregular coughing is usually caused by a respiratory tract infection but can also be triggered by choking, smoking, air pollution, asthma, gastroesophageal reflux disease, post-nasal drip, chronic bronchitis, lung tumors, heart failure and medications such as angiotensin-converting-enzyme inhibitors (ACE inhibitors) and beta blockers.
Treatment should target the cause; for example, smoking cessation or discontinuing ACE inhibitors. Cough suppressants such as codeine or dextromethorphan are frequently prescribed, but are not recommended for children. Other treatment options may target airway inflammation or may promote mucus expectoration. As it is a natural protective reflex, suppressing the cough reflex might have damaging effects, especially if the cough is productive (producing phlegm).
Presentation
Complications
The complications of coughing can be classified as either acute or chronic. Acute complications include cough syncope (fainting spells due to decreased blood flow to the brain when coughs are prolonged and forceful), insomnia, cough-induced vomiting, subconjunctival hemorrhage or "red eye", coughing defecation and in women with a prolapsed uterus, cough urination. Chronic complications are common and include abdominal or pelvic hernias, fatigue fractures of lower ribs and costochondritis. Chronic or violent coughing can contribute to damage to the pelvic floor and a possible cystocele.
Differential diagnosis
A cough in children may be either a normal physiological reflex or due to an underlying cause. In healthy children it may be normal in the absence of any disease to cough ten times a day. The most common cause of an acute or subacute cough is a viral respiratory tract infection. A healthy adult also coughs 18.6 times a day on average, but in the population with respiratory disease the geometric mean frequency is 275 times a day. In adults with a chronic cough, i.e. a cough longer than 8 weeks, more than 90% of cases are due to post-nasal drip, asthma, eosinophilic bronchitis, and gastroesophageal reflux disease. The causes of chronic cough are similar in children with the addition of bacterial bronchitis.
Infections
A cough can be the result of a respiratory tract infection such as the common cold, COVID-19, acute bronchitis, pneumonia, pertussis, or tuberculosis. In the vast majority of cases, acute coughs, i.e. coughs shorter than 3 weeks, are due to the common cold. In people with a normal chest X-ray, tuberculosis is a rare finding. Pertussis is increasingly being recognised as a cause of troublesome coughing in adults.
After a respiratory tract infection has cleared, the person may be left with a postinfectious cough. This typically is a dry, non-productive cough that produces no phlegm. Symptoms may include a tightness in the chest, and a tickle in the throat. This cough may often persist for weeks after an illness. The cause of the cough may be inflammation similar to that observed in repetitive stress disorders such as carpal tunnel syndrome. The repetition of coughing produces inflammation which produces discomfort, which in turn produces more coughing. Postinfectious cough typically does not respond to conventional cough treatments. Medication used for postinfectious coughs may include ipratropium to treat the inflammation, as well as cough suppressants to reduce frequency of the cough until inflammation clears. Inflammation may increase sensitivity to other existing issues such as allergies, and treatment of other causes of coughs (such as use of an air purifier or allergy medicines) may help speed recovery.
Reactive airway disease
When coughing is the only complaint of a person who meets the criteria for asthma (bronchial hyperresponsiveness and reversibility), this is termed cough-variant asthma. Atopic cough and eosinophilic bronchitis are related conditions. Atopic cough occurs in individuals with a family history of atopy (an allergic condition), abundant eosinophils in the sputum, but with normal airway function and responsiveness. Eosinophilic bronchitis is characterized by eosinophils in sputum and in bronchoalveolar lavage fluid without airway hyperresponsiveness or an atopic background. This condition responds to treatment with corticosteroids. Cough can also worsen in an acute exacerbation of chronic obstructive pulmonary disease.
Asthma is a common cause of chronic cough in adults and children. Coughing may be the only symptom the person has from their asthma, or asthma symptoms may also include wheezing, shortness of breath, and a tight feeling in their chest. Depending on how severe the asthma is, it can be treated with bronchodilators (medicine which causes the airways to open up) or inhaled steroids. Treatment of the asthma should make the cough go away.
Chronic bronchitis is defined clinically as a persistent cough that produces sputum (phlegm) and mucus, for at least three months in two consecutive years. Chronic bronchitis is often the cause of "smoker's cough". The tobacco smoke causes inflammation, secretion of mucus into the airway, and difficulty clearing that mucus out of the airways. Coughing helps clear those secretions out. May be treated by quitting smoking. May also be caused by pneumoconiosis and long-term fume inhalation.
Gastroesophageal reflux
In people with unexplained cough, gastroesophageal reflux disease should be considered. This occurs when acidic contents of the stomach come back up into the esophagus. Symptoms usually associated with GERD include heartburn, sour taste in the mouth, or a feeling of acid reflux in the chest, although, more than half of the people with cough from GERD do not have any other symptoms. An esophageal pH monitor can confirm the diagnosis of GERD. Sometimes GERD can complicate respiratory ailments related to cough, such as asthma or bronchitis. The treatment involves anti-acid medications and lifestyle changes with surgery indicated in cases not manageable with conservative measures.
Air pollution
Coughing may be caused by air pollution including tobacco smoke, particulate matter, irritant gases, and dampness in a home.
The human health effects of poor air quality are far reaching, but principally affect the body's respiratory system and the cardiovascular system. Individual reactions to air pollutants depend on the type of pollutant a person is exposed to, the degree of exposure, the individual's health status and genetics. People who exercise outdoors on hot, smoggy days, for example, increase their exposure to pollutants in the air.
Foreign body
A foreign body can sometimes be suspected, for example if the cough started suddenly when the patient was eating. Rarely, sutures left behind inside the airway branches can cause coughing. A cough can be triggered by dryness from mouth breathing or recurrent aspiration of food into the windpipe in people with swallowing difficulties.
Drug-induced cough
Drugs used for treatments other than coughs, such as ACE inhibitors which are often used to treat high blood pressure, can sometimes cause cough as a side effect, and stopping their use will stop the cough. Beta blockers similarly cause cough as an adverse event.
Habit cough
A habit cough is one that responds to behavioral or psychiatric therapy after organic causes have been excluded. Absence of the cough during sleep is common, but not diagnostic. A tic cough is thought to be more common in children than in adults.
Neurogenic cough
Some cases of chronic cough may be attributed to a sensory neuropathic disorder. Treatment for neurogenic cough may include the use of certain neuralgia medications. Coughing may occur in tic disorders such as Tourette syndrome, although it should be distinguished from throat-clearing in this disorder.
Other
Cough may also be caused by conditions affecting the lung tissue such as bronchiectasis, cystic fibrosis, interstitial lung diseases and sarcoidosis. Coughing can also be triggered by benign or malignant lung tumors or mediastinal masses. Through irritation of the nerve, diseases of the external auditory canal (wax, for example) can also cause cough. Cardiovascular diseases associated with cough are heart failure, pulmonary infarction and aortic aneurysm. Nocturnal cough is associated with heart failure, as the left ventricle doesn't effectively pump blood forward, resulting in blood being backed up in the pulmonary veins, which in turn causing pulmonary edema and resultant cough. Other causes of nocturnal cough include asthma, post-nasal drip and gastroesophageal reflux disease (GERD). Another cause of cough occurring preferentially in supine position is recurrent aspiration.
Given its irritant nature to mammal tissues, capsaicin is widely used to determine the cough threshold and as a tussive stimulant in clinical research of cough suppressants. Capsaicin is what makes chili peppers spicy, and might explain why workers in factories with these fruits can develop a cough.
Coughing may also be used for social reasons, and as such is not always involuntary. A voluntary cough, often written as "ahem", can be used to attract attention or express displeasure, as a form of nonverbal, paralingual metacommunication.
Airway clearance
Coughing, and huffing are important ways of removing mucus as sputum in many conditions such as cystic fibrosis, and chronic bronchitis.
Pathophysiology
A cough is a protective reflex in healthy individuals which is influenced by psychological factors. The cough reflex is initiated by stimulation of two different classes of afferent nerves, namely the myelinated rapidly adapting receptors, and nonmyelinated C-fibers with endings in the lung.
Diagnostic approach
The type of cough may help in the diagnosis. For instance, an inspiratory "whooping" sound on coughing almost doubles the likelihood that the illness is pertussis.
Blood may occur in small amounts with severe cough of many causes, but larger amounts suggests bronchitis, bronchiectasis, tuberculosis, or primary lung cancer.
Further workup may include labs, x-rays, and spirometry.
Classification
A cough can be classified by its duration, character, quality, and timing. The duration can be either acute (of sudden onset) if it is present less than three weeks, subacute if it is present between three or eight weeks, and chronic when lasting longer than eight weeks. A cough can be non-productive (dry) or productive (when phlegm is produced that may be coughed up as sputum). It may occur only at night (then called nocturnal cough), during both night and day, or just during the day.
A number of characteristic coughs exist. While these have not been found to be diagnostically useful in adults, they are of use in children. A barky cough is part of the common presentation of croup. A staccato cough has been classically described with neonatal chlamydial pneumonia.
Treatment
The treatment of a cough in children is based on the underlying cause. In children half of cases go away without treatment in 10 days and 90% in 25 days.
According to the American Academy of Pediatrics the use of cough medicine to relieve cough symptoms is supported by little evidence and thus not recommended for treating cough symptoms in children. There is tentative evidence that the use of honey is better than no treatment or diphenhydramine in decreasing coughing. It does not alleviate coughing to the same extent as dextromethorphan but it shortens the cough duration better than placebo and salbutamol. A trial of antibiotics or inhaled corticosteroids may be tried in children with a chronic cough in an attempt to treat protracted bacterial bronchitis or asthma respectively. There is insufficient evidence to recommend treating children who have a cough that is not related to a specific condition with inhaled anti-cholinergics.
Because coughing can spread disease through infectious aerosol droplets, it is recommended to cover one's mouth and nose with the forearm, the inside of the elbow, a tissue or a handkerchief while coughing.
Traditional medicine
One of the pharmaceutical dosage forms in traditional medicine for treatment of coughs was linctus. A linctus is a medicine in the form of a syrup, taken to relieve coughs and sore throats. The linctus is a syrup that helps relieve dry coughs.
Epidemiology
A cough is the most common reason for visiting a primary care physician in the United States.
Other animals
Marine mammals such as dolphins and whales cannot cough. Some invertebrates such as insects and spiders cannot cough or sneeze. Alligators can cough. Domestic animals and vertebrates such as dogs and cats can cough, because of diseases, allergies, dust or choking. In particular, cats are known for coughing before spitting up a hairball.
In other domestic animals, horses can cough because of infections, or due to poor ventilation and dust in enclosed spaces. Kennel cough in dogs can result from a viral or bacterial infection.
Deer can cough similarly to humans as a result of respiratory tract infections, such as parasitic bronchitis caused by a species of Dictyocaulus.
| Biology and health sciences | Symptoms and signs | Health |
197002 | https://en.wikipedia.org/wiki/Endoscopy | Endoscopy | An endoscopy is a procedure used in medicine to look inside the body. The endoscopy procedure uses an endoscope to examine the interior of a hollow organ or cavity of the body. Unlike many other medical imaging techniques, endoscopes are inserted directly into the organ.
There are many types of endoscopies. Depending on the site in the body and type of procedure, an endoscopy may be performed by either a doctor or a surgeon. A patient may be fully conscious or anaesthetised during the procedure. Most often, the term endoscopy is used to refer to an examination of the upper part of the gastrointestinal tract, known as an esophagogastroduodenoscopy.
For nonmedical use, similar instruments are called borescopes.
History
Adolf Kussmaul was fascinated by sword swallowers who would insert a sword down their throat without gagging. This drew inspiration to insert a hollow tube for observation; the next problem to solve was how to shine light through the tube, as they were still relying on candles and oil lamps as light sources.
The term endoscope was first used on February 7, 1855, by engineer-optician Charles Chevalier, in reference to the uréthroscope of Désormeaux, who himself began using the former term a month later. The self-illuminated endoscope was developed at Glasgow Royal Infirmary in Scotland (one of the first hospitals to have mains electricity) in 1894/5 by John Macintyre as part of his specialization in the investigation of the larynx.
Medical uses
Endoscopy may be used to investigate symptoms in the digestive system including nausea, vomiting, abdominal pain, difficulty swallowing, and gastrointestinal bleeding. It is also used in diagnosis, most commonly by performing a biopsy to check for conditions such as anemia, bleeding, inflammation, and cancers of the digestive system. The procedure may also be used for treatment such as cauterization of a bleeding vessel, widening a narrow esophagus, clipping off a polyp or removing a foreign object.
Specialty professional organizations that specialize in digestive problems advise that many patients with Barrett's esophagus receive endoscopies too frequently. Such societies recommend that patients with Barrett's esophagus and no cancer symptoms after two biopsies receive biopsies as indicated and no more often than the recommended rate.
Applications
Health care providers can use endoscopy to review any of the following body parts:
The gastrointestinal tract (GI tract):
oesophagus, stomach and duodenum (esophagogastroduodenoscopy)
small intestine (enteroscopy)
large intestine/colon (colonoscopy, sigmoidoscopy)
Magnification endoscopy
bile duct
endoscopic retrograde cholangiopancreatography (ERCP), duodenoscope-assisted cholangiopancreatoscopy, intraoperative cholangioscopy
rectum (rectoscopy) and anus (anoscopy), both also referred to as (proctoscopy)
The respiratory tract
The nose (rhinoscopy)
The upper respiratory tract (laryngoscopy)
The lower respiratory tract (bronchoscopy)
The ear (otoscope)
The Urinary bladder (cystoscopy)
The Ureter (ureteroscopy)
The female reproductive system (gynoscopy)
The cervix (colposcopy)
The uterus (hysteroscopy)
The fallopian tubes (falloposcopy)
Normally closed body cavities (through a small incision):
The abdominal or pelvic cavity (laparoscopy)
The interior of a joint (arthroscopy)
Organs of the chest (thoracoscopy and mediastinoscopy)
Endoscopy is used for many procedures:
During pregnancy
The amnion (amnioscopy)
The fetus (fetoscopy)
Plastic surgery
Panendoscopy (or triple endoscopy)
Combines laryngoscopy, esophagoscopy, and bronchoscopy
Orthopedic surgery
Hand surgery, such as endoscopic carpal tunnel release
Knee surgery, such as anterior cruciate ligament reconstruction
Epidural space (epiduroscopy)
Bursae (bursectomy)
Endodontic surgery
Maxillary sinus surgery
Apicoectomy
Endoscopic endonasal surgery
Endoscopic spinal surgery
Endoscopic nerve decompression for peripheral nerves
An endoscopy is a simple procedure that allows a doctor to look inside human bodies using an instrument called an endoscope. A cutting tool can be attached to the end of the endoscope, and the apparatus can then be used to perform minor procedures such as tissue biopsies, banding of oesophageal varices or removal of polyps.
Application in other fields
For non-medical use, such as internal inspection of complex technical systems, borescopes are used. These are similar to endoscopes.
The planning and architectural community use architectural endoscopy for pre-visualization of scale models of proposed buildings and cities
Endoscopes are also a tool helpful in the examination of improvised explosive devices by bomb disposal personnel.
Law enforcement uses endoscopes for conducting surveillance via tight spaces.
Risks
The main risks are infection, over-sedation, perforation, or a tear of the stomach or esophagus lining and bleeding. Although perforation generally requires surgery, certain cases may be treated with antibiotics and intravenous fluids. Bleeding may occur at the site of a biopsy or polyp removal. Such typically minor bleeding may simply stop on its own or be controlled by cauterisation. Seldom does surgery become necessary. Perforation and bleeding are rare during gastroscopy. Other minor risks include drug reactions and complications related to other diseases the patient may have. Consequently, patients should inform their doctor of all allergic tendencies and medical problems. Occasionally, the site of the sedative injection may become inflamed and tender for a short time. This is usually not serious and warm compresses for a few days are usually helpful. While any of these complications may possibly occur, each of them occurs quite infrequently. A doctor can further discuss risks with the patient with regard to the particular need for gastroscopy.
After the endoscopy
After the procedure, the patient will be observed and monitored by a qualified individual in the endoscopy room, or a recovery area, until a significant portion of the medication has worn off. Occasionally the patient is left with a mild sore throat, which may respond to saline gargles, or chamomile tea. It may last for weeks or not happen at all. The patient may have a feeling of distention from the insufflated air that was used during the procedure. Both problems are mild and fleeting. When fully recovered, the patient will be instructed when to resume their usual diet (probably within a few hours) and will be allowed to be taken home. Where sedation has been used, most facilities mandate that the patient be taken home by another person and that they not drive or handle machinery for the remainder of the day. Patients who have had an endoscopy without sedation are able to leave unassisted.
| Technology | Diagnostic technologies | null |
197020 | https://en.wikipedia.org/wiki/Gallbladder | Gallbladder | In vertebrates, the gallbladder, also known as the cholecyst, is a small hollow organ where bile is stored and concentrated before it is released into the small intestine. In humans, the pear-shaped gallbladder lies beneath the liver, although the structure and position of the gallbladder can vary significantly among animal species. It receives bile, produced by the liver, via the common hepatic duct, and stores it. The bile is then released via the common bile duct into the duodenum, where the bile helps in the digestion of fats.
The gallbladder can be affected by gallstones, formed by material that cannot be dissolved – usually cholesterol or bilirubin, a product of hemoglobin breakdown. These may cause significant pain, particularly in the upper-right corner of the abdomen, and are often treated with removal of the gallbladder (called a cholecystectomy). Cholecystitis, inflammation of the gallbladder, has a wide range of causes, including result from the impaction of gallstones, infection, and autoimmune disease.
Structure
The human gallbladder is a hollow grey-blue organ that sits in a shallow depression below the right lobe of the liver. In adults, the gallbladder measures approximately in length and in diameter when fully distended. The gallbladder has a capacity of about .
The gallbladder is shaped like a pear, with its tip opening into the cystic duct. The gallbladder is divided into three sections: the fundus, body, and neck. The fundus is the rounded base, angled so that it faces the abdominal wall. The body lies in a depression in the surface of the lower liver. The neck tapers and is continuous with the cystic duct, part of the biliary tree. The gallbladder fossa, against which the fundus and body of the gallbladder lie, is found beneath the junction of hepatic segments IVB and V. The cystic duct unites with the common hepatic duct to become the common bile duct. At the junction of the neck of the gallbladder and the cystic duct, there is an out-pouching of the gallbladder wall forming a mucosal fold known as "Hartmann's pouch".
Lymphatic drainage of the gallbladder follows the cystic node, which is located between the cystic duct and the common hepatic duct. Lymphatics from the lower part of the organ drain into lower hepatic lymph nodes. All the lymph finally drains into celiac lymph nodes.
Microanatomy
The gallbladder wall is composed of a number of layers. The innermost surface of the gallbladder wall is lined by a single layer of columnar cells with a brush border of microvilli, very similar to intestinal absorptive cells. Underneath the epithelium is an underlying lamina propria, a muscular layer, an outer perimuscular layer and serosa. Unlike elsewhere in the intestinal tract, the gallbladder does not have a muscularis mucosae, and the muscular fibres are not arranged in distinct layers.
The mucosa, the inner portion of the gallbladder wall, consists of a lining of a single layer of columnar cells, with cells possessing small hair-like attachments called microvilli. This sits on a thin layer of connective tissue, the lamina propria. The mucosa is curved and collected into tiny outpouchings called rugae.
A muscular layer sits beneath the mucosa. This is formed by smooth muscle, with fibres that lie in longitudinal, oblique and transverse directions, and are not arranged in separate layers. The muscle fibres here contract to expel bile from the gallbladder. A distinctive feature of the gallbladder is the presence of Rokitansky–Aschoff sinuses, deep outpouchings of the mucosa that can extend through the muscular layer, and which indicate adenomyomatosis. The muscular layer is surrounded by a layer of connective and fat tissue.
The outer layer of the fundus of gallbladder, and the surfaces not in contact with the liver, are covered by a thick serosa, which is exposed to the peritoneum. The serosa contains blood vessels and lymphatics. The surfaces in contact with the liver are covered in connective tissue.
Variation
The gallbladder varies in size, shape, and position among different people. Rarely, two or even three gallbladders may coexist, either as separate bladders draining into the cystic duct, or sharing a common branch that drains into the cystic duct. Additionally, the gallbladder may fail to form at all. Gallbladders with two lobes separated by a septum may also exist. These abnormalities are not likely to affect function and are generally asymptomatic.
The location of the gallbladder in relation to the liver may also vary, with documented variants including gallbladders found within, above, on the left side of, behind, and detached or suspended from the liver. Such variants are very rare: from 1886 to 1998, only 110 cases of left-lying liver, or less than one per year, were reported in scientific literature.
An anatomical variation can occur, known as a Phrygian cap, which is an innocuous fold in the fundus, named after its resemblance to the Phrygian cap.
Development
The gallbladder develops from an endodermal outpouching of the embryonic gut tube. Early in development, the human embryo has three germ layers and abuts an embryonic yolk sac. During the second week of embryogenesis, as the embryo grows, it begins to surround and envelop portions of this sac. The enveloped portions form the basis for the adult gastrointestinal tract. Sections of this foregut begin to differentiate into the organs of the gastrointestinal tract, such as the esophagus, stomach, and intestines.
During the fourth week of embryological development, the stomach rotates. The stomach, originally lying in the midline of the embryo, rotates so that its body is on the left. This rotation also affects the part of the gastrointestinal tube immediately below the stomach, which will go on to become the duodenum. By the end of the fourth week, the developing duodenum begins to spout a small outpouching on its right side, the hepatic diverticulum, which will go on to become the biliary tree. Just below this is a second outpouching, known as the cystic diverticulum, that will eventually develop into the gallbladder.
Function
The main functions of the gallbladder are to store and concentrate bile, also called gall, needed for the digestion of fats in food. Produced by the liver, bile flows through small vessels into the larger hepatic ducts and ultimately through the cystic duct (parts of the biliary tree) into the gallbladder, where it is stored. At any one time, of bile is stored within the gallbladder.
When food containing fat enters the digestive tract, it stimulates the secretion of cholecystokinin (CCK) from I cells of the duodenum and jejunum. In response to cholecystokinin, the gallbladder rhythmically contracts and releases its contents into the common bile duct, eventually draining into the duodenum. The bile emulsifies fats in partly digested food, thereby assisting their absorption. Bile consists primarily of water and bile salts, and also acts as a means of eliminating bilirubin, a product of hemoglobin metabolism, from the body.
The bile that is secreted by the liver and stored in the gallbladder is not the same as the bile that is secreted by the gallbladder. During gallbladder storage of bile, it is concentrated 3–10 fold by removal of some water and electrolytes. This is through the active transport of sodium and chloride ions across the epithelium of the gallbladder, which creates an osmotic pressure that also causes water and other electrolytes to be reabsorbed.
A function of the gallbladder appears to be protection against carcinogenesis as indicated by observations that removal of the gallbladder (cholecystectomy) increases subsequent cancer risk. For instance, a systematic review and meta analysis of eighteen studies concluded that cholecystectomy has a harmful effect on the risk of right-sided colon cancer. Another recent study reported a significantly increased total cancer risk, including increased risk of several different types of cancer, after cholecystectomy.
Clinical significance
Gallstones
Gallstones form when the bile is saturated, usually with either cholesterol or bilirubin. Most gallstones do not cause symptoms, with stones either remaining in the gallbladder or passed along the biliary system. When symptoms occur, severe "colicky" pain in the upper right quadrant of the abdomen is often felt. If the stone blocks the gallbladder, inflammation known as cholecystitis may result. If the stone lodges in the biliary system, jaundice may occur; if the stone blocks the pancreatic duct, pancreatitis may occur. Gallstones are diagnosed using ultrasound. When a symptomatic gallstone occurs, it is often managed by waiting for it to be passed naturally. Given the likelihood of recurrent gallstones, surgery to remove the gallbladder is often considered. Some medication, such as ursodeoxycholic acid, may be used; lithotripsy, a non-invasive mechanical procedure used to break down the stones, may also be used.
Inflammation
Known as cholecystitis, inflammation of the gallbladder is commonly caused by obstruction of the duct with gallstones, which is known as cholelithiasis. Blocked bile accumulates, and pressure on the gallbladder wall may lead to the release of substances that cause inflammation, such as phospholipase. There is also the risk of bacterial infection. An inflamed gallbladder is likely to cause sharp and localised pain, fever, and tenderness in the upper, right corner of the abdomen, and may have a positive Murphy's sign. Cholecystitis is often managed with rest and antibiotics, particularly cephalosporins and, in severe cases, metronidazole. Additionally the gallbladder may need to be removed surgically if inflammation has progressed far enough.
Gallbladder removal
A cholecystectomy is a procedure in which the gallbladder is removed. It may be removed because of recurrent gallstones and is considered an elective procedure. A cholecystectomy may be an open procedure, or a laparoscopic one. In the surgery, the gallbladder is removed from the neck to the fundus, and so bile will drain directly from the liver into the biliary tree. About 30 percent of patients may experience some degree of indigestion following the procedure, although severe complications are much rarer. About 10 percent of surgeries lead to a chronic condition of postcholecystectomy syndrome.
Complication
Biliary injury (bile duct injury) is the traumatic damage of the bile ducts. It is most commonly an iatrogenic complication of cholecystectomy — surgical removal of gall bladder, but can also be caused by other operations or by major trauma. The risk of biliary injury is more during laparoscopic cholecystectomy than during open cholecystectomy. Biliary injury may lead to several complications and may even cause death if not diagnosed in time and managed properly. Ideally biliary injury should be managed at a center with facilities and expertise in endoscopy, radiology and surgery.
Biloma is collection of bile within the abdominal cavity. It happens when there is a bile leak, for example after surgery for removing the gallbladder (laparoscopic cholecystectomy), with an incidence of 0.3–2%. Other causes are biliary surgery, liver biopsy, abdominal trauma, and, rarely, spontaneous perforation.
Cancer
Cancer of the gallbladder is uncommon and mostly occurs in later life. When cancer occurs, it is mostly of the glands lining the surface of the gallbladder (adenocarcinoma). Gallstones are thought to be linked to the formation of cancer. Other risk factors include large (>1 cm) gallbladder polyps and having a highly calcified "porcelain" gallbladder.
Cancer of the gallbladder can cause attacks of biliary pain, yellowing of the skin (jaundice), and weight loss. A large gallbladder may be able to be felt in the abdomen. Liver function tests may be elevated, particularly involving GGT and ALP, with ultrasound and CT scans being considered medical imaging investigations of choice. Cancer of the gallbladder is managed by removing the gallbladder, however, the prognosis remains poor.
Cancer of the gallbladder may also be found incidentally after surgical removal of the gallbladder, with 1–3% of cancers identified in this way. Gallbladder polyps are mostly benign growths or lesions resembling growths that form in the gallbladder wall, and are only associated with cancer when they are larger in size (>1 cm). Cholesterol polyps, often associated with cholesterolosis ("strawberry gallbladder", a change in the gallbladder wall due to excess cholesterol), often cause no symptoms and are thus often detected in this way.
Tests
Tests used to investigate for gallbladder disease include blood tests and medical imaging. A full blood count may reveal an increased white cell count suggestive of inflammation or infection. Tests such as bilirubin and liver function tests may reveal if there is inflammation linked to the biliary tree or gallbladder, and whether this is associated with inflammation of the liver, and a lipase or amylase may be elevated if there is pancreatitis. Bilirubin may rise when there is obstruction of the flow of bile. A CA 19-9 level may be taken to investigate for cholangiocarcinoma.
An ultrasound is often the first medical imaging test performed when gallbladder disease such as gallstones are suspected. An abdominal X-ray or CT scan is another form of imaging that may be used to examine the gallbladder and surrounding organs. Other imaging options include MRCP (magnetic resonance cholangiopancreatography), ERCP and percutaneous or intraoperative cholangiography. A cholescintigraphy scan is a nuclear imaging procedure used to assess the condition of the gallbladder.
Other animals
Most vertebrates have gallbladders, but the form and arrangement of the bile ducts may vary considerably. In many species, for example, there are several separate ducts running to the intestine, rather than the single common bile duct found in humans. Several species of mammals (including horses, deer, rats, and laminoids), several species of birds (such as pigeons and some psittacine species), lampreys and all invertebrates do not have a gallbladder.
The bile from several species of bears is used in traditional Chinese medicine; bile bears are kept alive in captivity while their bile is extracted, in an industry characterized by animal cruelty.
History
Depictions of the gallbladder and biliary tree are found in Babylonian models found from 2000 BCE, and in ancient Etruscan model from 200 BCE, with models associated with divine worship.
Diseases of the gallbladder are known to have existed in humans since antiquity, with gallstones found in the mummy of Princess Amenen of Thebes dating to 1500 BCE. Some historians believe the death of Alexander the Great may have been associated with an acute episode of cholecystitis. The existence of the gallbladder has been noted since the 5th century, but it is only relatively recently that the function and the diseases of the gallbladder has been documented, particularly in the last two centuries.
The first descriptions of gallstones appear to have been in the Renaissance, perhaps because of the low incidence of gallstones in earlier times owing to a diet with more cereals and vegetables and less meat. Anthonius Benevinius in 1506 was the first to draw a connection between symptoms and the presence of gallstones. Ludwig Georg Courvoisier, after examining a number of cases in 1890 that gave rise to the eponymous Courvoisier's law, stated that in an enlarged, nontender gallbladder, the cause of jaundice is unlikely to be gallstones.
The first surgical removal of a gallstone (cholecystolithotomy) was in 1676 by physician Joenisius, who removed the stones from a spontaneously occurring biliary fistula. Stough Hobbs in 1867 performed the first recorded cholecystotomy, although such an operation was in fact described earlier by French surgeon Jean Louis Petit in the mid eighteenth century. German surgeon Carl Langenbuch performed the first cholecystectomy in 1882 for a sufferer of cholelithiasis. Before this, surgery had focused on creating a fistula for drainage of gallstones. Langenbuch reasoned that given several other species of mammal have no gallbladder, humans could survive without one.
The debate whether surgical removal of the gallbladder or simply gallstones was preferred was settled in the 1920s, with the consensus that removal of the gallbladder was preferred. It was only in the mid and late parts of the twentieth century that medical imaging techniques such as use of contrast medium and CT scans were used to view the gallbladder. The first laparoscopic cholecystectomy performed by Erich Mühe of Germany in 1985, although French surgeons Phillipe Mouret and Francois Dubois are often credited for their operations in 1987 and 1988 respectively.
Society and culture
To have "gall" is associated with bold, belligerent behaviour, whereas to have "bile" is associated with sourness.
In the Chinese medicine, the gallbladder () is associated with the Wuxing element of wood, in excess its emotion is belligerence and in deficiency cowardice and judgement, in the Chinese language it is related to a myriad of idioms, including using terms such as "a body completely [of] gall" () to describe a forward person, and "single, alone gallbladder hero" () to describe a lone hero, or "they have a lot of gall to talk like that".
In the Zangfu theory of Chinese medicine it is an extraordinary Fu or yang organ, as it holds bile. The gallbladder not only has a digestive role, but is seen as the seat of decision-making and judgement.
| Biology and health sciences | Gastrointestinal tract | Biology |
197037 | https://en.wikipedia.org/wiki/Artificial%20organ | Artificial organ | An artificial organ is a human-made organ device or tissue that is implanted or integrated into a humaninterfacing with living tissueto replace a natural organ, to duplicate or augment a specific function or functions so the patient may return to a normal life as soon as possible. The replaced function does not have to be related to life support, but it often is. For example, replacement bones and joints, such as those found in hip replacements, could also be considered artificial organs.
Implied by definition, is that the device must not be continuously tethered to a stationary power supply or other stationary resources such as filters or chemical processing units. (Periodic rapid recharging of batteries, refilling of chemicals, and/or cleaning/replacing of filters would exclude a device from being called an artificial organ.) Thus, a dialysis machine, while a very successful and critically important life support device that almost completely replaces the duties of a kidney, is not an artificial organ.
Purpose
Constructing and installing artificial organs, an extremely research-intensive and expensive process initially, may entail many years of ongoing maintenance services not needed by a natural organ:
providing life support to prevent imminent death while awaiting a transplant (e.g. artificial heart);
dramatically improving the patient's ability for self care (e.g. artificial limb);
improving the patient's ability to interact socially (e.g. cochlear implant); or
improving a patient's quality of life through cosmetic restoration after cancer surgery or an accident.
The use of any artificial organ by humans is almost always preceded by extensive experiments with animals. Initial testing in humans is frequently limited to those either already facing death or who have exhausted every other treatment possibility.
Examples
Artificial limbs
Artificial arms and legs, or prosthetics, are intended to restore a degree of normal function to amputees. Mechanical devices that allow amputees to walk again or continue to use two hands have probably been in use since ancient times, the most notable one being the simple peg leg. Since then, the development of artificial limbs has progressed rapidly. New plastics and other materials, such as carbon fiber have allowed artificial limbs to become stronger and lighter, limiting the amount of extra energy necessary to operate the limb. Additional materials have allowed artificial limbs to look much more realistic. Prostheses can roughly be categorized as upper- and lower-extremity and can take many shapes and sizes.
New advances in artificial limbs include additional levels of integration with the human body. Electrodes can be placed into nervous tissue, and the body can be trained to control the prosthesis. This technology has been used in both animals and humans. The prosthetic can be controlled by the brain using a direct implant or implant into various muscles.
Bladder
The two main methods for replacing bladder function involve either redirecting urine flow or replacing the bladder in situ. Standard methods for replacing the bladder involve fashioning a bladder-like pouch from intestinal tissue. As of 2017 methods to grow bladders using stem cells had been attempted in clinical research but this procedure was not part of medicine.
Brain
Neural prostheses are a series of devices that can substitute a motor, sensory or cognitive modality that might have been damaged as a result of an injury or a disease.
Neurostimulators, including deep brain stimulators, send electrical impulses to the brain in order to treat neurological and movement disorders, including Parkinson's disease, epilepsy, treatment resistant depression, and other conditions such as urinary incontinence. Rather than replacing existing neural networks to restore function, these devices often serve by disrupting the output of existing malfunctioning nerve centers to eliminate symptoms.
Scientists in 2013 created a mini brain that developed key neurological components until the early gestational stages of fetal maturation.
Corpora cavernosa
To treat erectile dysfunction, both corpora cavernosa can be irreversibly surgically replaced with manually inflatable penile implants. This is a drastic therapeutic surgery meant only for men who have complete impotence resistant to all other treatment approaches. An implanted pump in the groin or scrotum can be manipulated by hand to fill these artificial cylinders, normally sized to be direct replacements for the natural corpora cavernosa, from an implanted reservoir in order to achieve an erection.
Ear
In cases when a person is profoundly deaf or severely hard of hearing in both ears, a cochlear implant may be surgically implanted. Cochlear implants bypass most of the peripheral auditory system to provide a sense of sound via a microphone and some electronics that reside outside the skin, generally behind the ear. The external components transmit a signal to an array of electrodes placed in the cochlea, which in turn stimulates the cochlear nerve.
In the case of an outer ear trauma, a craniofacial prosthesis may be necessary.
Thomas Cervantes and his colleagues, who are from Massachusetts General Hospital, built an artificial ear from sheep cartilage by a 3D printer. With a lot of calculations and models, they managed to build an ear shaped like a typical human one. Modeled by a plastic surgeon, they had to adjust several times so the artificial ear can have curves and lines just like a human ear. The researchers said "The technology is now under development for clinical trials, and thus we have scaled up and redesigned the prominent features of the scaffold to match the size of an adult human ear and to preserve the aesthetic appearance after implantation." Their artificial ears have not been announced as successful, but they are still currently developing the project. Each year, thousands of children were born with a congenital deformity called microtia, where the external ear does not fully develop. This could be a major step forward in medical and surgical microtia treatment.
Eye
The most successful function-replacing artificial eye so far is actually an external miniature digital camera with a remote unidirectional electronic interface implanted on the retina, optic nerve, or other related locations inside the brain. The present state of the art yields only partial functionality, such as recognizing levels of brightness, swatches of color, and/or basic geometric shapes, proving the concept's potential.
Various researchers have demonstrated that the retina performs strategic image preprocessing for the brain. The problem of creating a completely functional artificial electronic eye is even more complex. Advances towards tackling the complexity of the artificial connection to the retina, optic nerve, or related brain areas, combined with ongoing advances in computer science, are expected to dramatically improve the performance of this technology.
Heart
Cardiovascular-related artificial organs are implanted in cases where the heart, its valves, or another part of the circulatory system is in disorder. The artificial heart is typically used to bridge the time to heart transplantation, or to permanently replace the heart in case heart transplantation is impossible. Artificial pacemakers represent another cardiovascular device that can be implanted to either intermittently augment (defibrillator mode), continuously augment, or completely bypass the natural living cardiac pacemaker as needed. Ventricular assist devices are another alternative, acting as mechanical circulatory devices that partially or completely replace the function of a failing heart, without the removal of the heart itself.
Besides these, lab-grown hearts and 3D bioprinted hearts are also being researched. Currently, scientists are limited in their ability to grow and print hearts due to difficulties in getting blood vessels and lab-made tissues to function cohesively.
Liver
HepaLife is developing a bioartificial liver device intended for the treatment of liver failure using stem cells. The artificial liver is designed to serve as a supportive device, either allowing the liver to regenerate upon failure, or to bridge the patient's liver functions until transplant is available. It is only made possible by the fact that it uses real liver cells (hepatocytes), and even then, it is not a permanent substitute.
Lungs
With some almost fully functional, artificial lungs promise to be a great success in the near future. An Ann Arbor company MC3 is currently working on this type of medical device.
Extracorporeal membrane oxygenation (ECMO) can be used to take significant load off of the native lung tissue and heart. In ECMO, one or more catheters are placed into the patient and a pump is used to flow blood over hollow membrane fibers, which exchange oxygen and carbon dioxide with the blood. Similar to ECMO, Extracorporeal Removal (ECCO2R) has a similar set-up, but mainly benefits the patient through carbon dioxide removal, rather than oxygenation, with the goal of allowing the lungs to relax and heal.
Ovaries
The ground work for the development of the artificial ovary was laid in the early 1990s.
Reproductive age patients who develop cancer often receive chemotherapy or radiation therapy, which damages oocytes and leads to early menopause. An artificial human ovary has been developed at Brown University with self-assembled microtissues created using novel 3-D petri dish technology. In a study funded and conducted by the NIH in 2017, scientists were successful in printing 3-D ovaries and implanting them in sterile mice. In the future, scientists hope to replicate this in larger animals as well as humans. The artificial ovary will be used for the purpose of in vitro maturation of immature oocytes and the development of a system to study the effect of environmental toxins on folliculogenesis.
Pancreas
An artificial pancreas is used to substitute endocrine functionality of a healthy pancreas for diabetic and other patients who require it. It can be used to improve insulin replacement therapy until glycemic control is practically normal as evident by the avoidance of the complications of hyperglycemia, and it can also ease the burden of therapy for the insulin-dependent. Approaches include using an insulin pump under closed loop control, developing a bio-artificial pancreas consisting of a biocompatible sheet of encapsulated beta cells, or using gene therapy.
Red blood cells
Artificial red blood cells (RBC) have already been in projects for about 60 years, but they started getting interest when the HIV-contaminated-donor blood crisis. Artificial RBCs will be dependent 100% on nanotechnology. A successful artificial RBC should be able to totally replace human RBC, which means it can carry on all the functions that a human RBC does.
The first artificial RBC, made by Chang and Poznanski in 1968, was made to transport Oxygen and Carbon Dioxide, also fulfilled antioxidant functions.
Scientists are working on a new kind of artificial RBC, which is one-fiftieth the size of a human RBC. They are made from purified human hemoglobin proteins that have been coated with a synthetic polymer. Thanks to the special materials of the artificial RBC, they can capture oxygen when blood pH is high, and release oxygen when blood pH is low. The polymer coating also keeps the hemoglobin from reacting with nitric oxide in the bloodstream, thus preventing dangerous constriction of the blood vessels. Allan Doctor, MD, stated that the artificial RBC can be used by anyone, with any blood type because the coating is immune silent.
Testes
Men whom have sustained testicular abnormalities through birth defects or injury have been able to replace the damaged testicle with a testicular prosthesis. Although the prosthetic does not restore biological reproductive function, the device has been shown to improve mental health for these patients.
Thymus
An implantable machine that performs the function of a thymus does not exist. However, researchers have been able to grow a thymus from reprogrammed fibroblasts. They expressed hope that the approach could one day replace or supplement neonatal thymus transplantation.
As of 2017, researchers at UCLA developed an artificial thymus that, although not yet implantable, is capable of performing all functions of a true thymus.
The artificial thymus would play an important role in the immune system, and it would use blood stem cells to produce more T cells, which in turn, help the body fight infections. It would ultimately give the body a better ability to fight cancer cells. As people age, if their thymus stops working well, an artificial thymus could also be a potentially viable option.
The idea of using T cells to fight against infections has been around for a time, but until recently, the idea of using a T cell source, an artificial thymus is proposed. "We know that the key to creating a consistent and safe supply of cancer-fighting T cells would be to control the process in a way that deactivates all T cell receptors in the transplanted cells, except for the cancer-fighting receptors," said Dr. Gay Crooks of UCLA. The scientist also found that the T cells produced by the artificial thymus carried a diverse range of T cell receptors and worked similarly to the T cells produced by a normal thymus. Since they can work like human thymus, artificial thymus can supply a consistent amount of T cells to the body for the patients who are in need of treatments.
Trachea
The field of artificial tracheas went through a period of high interest and excitement with the work of Paolo Macchiarini at the Karolinska Institute and elsewhere from 2008 to around 2014, with front-page coverage in newspapers and on television. Concerns were raised about his work in 2014 and by 2016 he had been fired and high level management at Karolinska had been dismissed, including people involved in the Nobel Prize.
As of 2017 engineering a tracheaa hollow tube lined with cellshad proved more challenging than originally thought; challenges include the difficult clinical situation of people who present as clinical candidates, who generally have been through multiple procedures already; creating an implant that can become fully developed and integrate with host while withstanding respiratory forces, as well as the rotational and longitudinal movement the trachea undergoes.
Enhancement
It is also possible to construct and install an artificial organ to give its possessor abilities that are not naturally occurring. Research is proceeding in areas of vision, memory, and information processing. Some current research focuses on restoring short-term memory in accident victims and long-term memory in dementia patients.
One area of success was achieved when Kevin Warwick carried out a series of experiments extending his nervous system over the internet to control a robotic hand and the first direct electronic communication between the nervous systems of two humans.
This might also include the existing practice of implanting subcutaneous chips for identification and location purposes (ex. RFID tags).
Microchips
Organ chips are devices containing hollow microvessels filled with cells simulating tissue and/or organs as a microfluidic system that can provide key chemical and electrical signal information. This is distinct from an alternative use of the term microchip, which refers to small, electronic chips that are commonly used as an identifier and can also contain a transponder.
This information can create various applications such as creating "human in vitro models" for both healthy and diseased organs, drug advancements in toxicity screening as well as replacing animal testing.
Using 3D cell culture techniques enables scientists to recreate the complex extracellular matrix, ECM, found in in vivo to mimic human response to drugs and human diseases.
Organs on chips are used to reduce the failure rate in new drug development; microengineering these allows for a microenvironment to be modeled as an organ.
| Technology | Devices | null |
197112 | https://en.wikipedia.org/wiki/Bumblebee | Bumblebee | A bumblebee (or bumble bee, bumble-bee, or humble-bee) is any of over 250 species in the genus Bombus, part of Apidae, one of the bee families. This genus is the only extant group in the tribe Bombini, though a few extinct related genera (e.g., Calyptapis) are known from fossils. They are found primarily in higher altitudes or latitudes in the Northern Hemisphere, although they are also found in South America, where a few lowland tropical species have been identified. European bumblebees have also been introduced to New Zealand and Tasmania. Female bumblebees can sting repeatedly, but generally ignore humans and other animals.
Most bumblebees are social insects that form colonies with a single queen. The colonies are smaller than those of honey bees, growing to as few as 50 individuals in a nest. Cuckoo bumblebees are brood parasitic and do not make nests or form colonies; their queens aggressively invade the nests of other bumblebee species, kill the resident queens and then lay their own eggs, which are cared for by the resident workers. Cuckoo bumblebees were previously classified as a separate genus, but are now usually treated as members of Bombus.
Bumblebees have round bodies covered in soft hair (long branched setae) called 'pile', making them appear and feel fuzzy. They have aposematic (warning) coloration, often consisting of contrasting bands of colour, and different species of bumblebee in a region often resemble each other in mutually protective Müllerian mimicry. Harmless insects such as hoverflies often derive protection from resembling bumblebees, in Batesian mimicry, and may be confused with them. Nest-making bumblebees can be distinguished from similarly large, fuzzy cuckoo bumblebees by the form of the female hind leg. In nesting bumblebees, it is modified to form a pollen basket, a bare shiny area surrounded by a fringe of hairs used to transport pollen, whereas in cuckoo bumblebees, the hind leg is hairy all round, and they never carry pollen.
Like their relatives the honeybees, bumblebees feed on nectar, using their long hairy tongues to lap up the liquid; the proboscis is folded under the head during flight. Bumblebees gather nectar to add to the stores in the nest, and pollen to feed their young. They forage using colour and spatial relationships to identify flowers to feed from. Some bumblebees steal nectar, making a hole near the base of a flower to access the nectar while avoiding pollen transfer. Bumblebees are important agricultural pollinators, so their decline in Europe, North America, and Asia is a cause for concern. The decline has been caused by habitat loss, the mechanisation of agriculture, and pesticides.
Etymology and common names
The word "bumblebee" is a compound of "bumble" and "bee"—'bumble' meaning to hum, buzz, drone, or move ineptly or flounderingly. The generic name Bombus, assigned by Pierre André Latreille in 1802, is derived from the Latin word for a buzzing or humming sound, borrowed from Ancient Greek ().
According to the Oxford English Dictionary (OED), the term "bumblebee" was first recorded as having been used in the English language in the 1530 work Lesclarcissement by John Palsgrave, "I bomme, as a bombyll bee dothe." However the OED also states that the term "humblebee" predates it, having first been used in 1450 in Fysshynge wyth Angle, "In Juyll the greshop & the humbylbee in the medow." The latter term was used in A Midsummer Night's Dream () by William Shakespeare, "The honie-bags steale from the humble Bees." Similar terms are used in other Germanic languages, such as the German (Old High German ), Dutch or Swedish .
An old provincial name, "dumbledor", also denoted a buzzing insect such as a bumblebee or cockchafer, "dumble" probably imitating the sound of these insects, while "dor" meant "beetle".
In On the Origin of Species (1859), Charles Darwin speculated about "humble-bees" and their interactions with other species:
However, "bumblebee" remained in use, for example in The Tale of Mrs. Tittlemouse (1910) by Beatrix Potter, "Suddenly round a corner, she met Babbitty Bumble--"Zizz, Bizz, Bizzz!" said the bumble bee." Since World War II "humblebee" has fallen into near-total disuse.
Phylogeny
The bumblebee tribe Bombini is one of four groups of corbiculate bees (those with pollen baskets) in the Apidae, the others being the Apini (honey bees), Euglossini (orchid bees), and Meliponini (stingless bees). The corbiculate bees are a monophyletic group. Advanced eusocial behaviour appears to have evolved twice in the group, giving rise to controversy, now largely settled, as to the phylogenetic origins of the four tribes; it had been supposed that eusocial behaviour had evolved only once, requiring the Apini to be close to the Meliponini, which they do not resemble. It is now thought that the Apini (with advanced societies) and Euglossini are closely related, while the primitively eusocial Bombini are close to the Meliponini, which have somewhat more advanced eusocial behaviour. Sophie Cardinal and Bryan Danforth comment that "While remarkable, a hypothesis of dual origins of advanced eusociality is congruent with early studies on corbiculate morphology and social behavior." Their analysis, combining molecular, morphological and behavioural data, gives the following cladogram:
On this hypothesis, the molecular data suggest that the Bombini are 25 to 40 million years old, while the Meliponini (and thus the clade that includes the Bombini and Meliponini) are 81 to 96 million years old, about the same age as the corbiculate group.
However, a more recent phylogeny using transcriptome data from 3,647 genes of ten corbiculate bee species supports the single origin of eusociality hypothesis in the corbiculate bees. They find that Bombini is in fact sister to Meliponini, corroborating that previous finding from Sophie Cardinal and Bryan Danforth (2011). However, Romiguier et al. (2015) shows that Bombini, Meliponini, and Apini form a monophyletic group, where Apini shares a most recent common ancestor with the Bombini and Meliponini clade, while Euglossini is most distantly related to all three, since it does not share the same most recent common ancestor as Bombini, Meliponini, and Apini. Thus, their analysis supports the single origin of eusociality hypothesis within the corbiculate bees, where eusociality evolved in the common ancestor of Bombini, Apini, and Meliponini.
The fossil record for bees is limited, with around 14 species that might possibly be Bombini having been described by 2019. The only Bombus relatives in Bombini are the late Eocene Calyptapis florissantensis from the Florissant Formation, USA, and Oligobombus cuspidatus from the Bembridge Marls of the Isle of Wight. Two species of Bombus have been described from the Oligocene of Beşkonak, Bucak Turkey: Bombus (Mendacibombus) beskonakensis and Bombus (Paraelectrobombus) patriciae. Both species were originally placed in genera considered at the time of description as outside of Bombus, being initially named Oligoapis beskonakensis and Paraelectrobombus patriciae respectively, however reexaminiation of the fore-wings lead to both being considered as Bombus species In 2012 a fossil bumblebee from the Miocene was found in Germany's Randeck Maar and classified as Bombus (Bombus) randeckensis. In 2014, another species, Bombus cerdanyensis, was described from Late Miocene lacustrine beds of La Cerdanya, Spain, but not initially placed into any subgenus, The species Bombus trophonius was described in October 2017 and placed in Bombus subgenus Cullumanobombus. A redescription of the Bombini fossil record by Dehon et al (2019) resulted in the synonymization of the genus Oligoapis with Bombus subgenus Mendacibombus, and the placement of genus Paraelectrobombus as Bombus subgenus Paraelectrobombus, rather than as a genus in Electrobombini. The subgenus Cullumanobombus was expanded to include not only Bombus trophonius but also Bombus randeckensis which was moved from subgenus Bombus and Bombus pristinus, first described by Unger (1867). Within the subgenus Melanobombus only Bombus cerdanyensis is present from the fossil record. An additional three species, "Bombus" luianus, "Bombus" anacolus and "Bombus" dilectus have been attributed to Bombus from the Middle Miocene Shanwang formation of China by Zhang, (1990) and Zhang et al (1994). Due to not being able to study Zhang's type specimens, but only illustrations of the fossils, Dehon et al did not place the three species within any specific subgenera, and considered all three as "species inquirenda", needing fuller re-examination. Two other species were not examined at all by Dehon et al, Bombus? crassipes of the Late Miocene Krottensee deposits in the Czech Republic, and Bombus proavus from the Middle Miocene Latah Formation, USA.
Taxonomy
The genus Bombus, the only one extant genus in the tribe Bombini, comprises over 250 species; for an overview of the differences between bumblebees and other bees and wasps, see characteristics of common wasps and bees. The genus has been divided variously into up to 49 subgenera, a degree of complexity criticised by Williams (2008). The cuckoo bumblebees Psithyrus have sometimes been treated as a separate genus but are now considered to be part of Bombus, in one or more subgenera.
Examples of Bombus species include Bombus pauloensis, Bombus dahlbomii, Bombus fervidus, Bombus lapidarius, Bombus ruderatus, and Bombus rupestris.
Subgenera of the genus Bombus
General description
Bumblebees vary in appearance, but are generally plump and densely furry. They are larger, broader and stouter-bodied than honeybees, and their abdomen tip is more rounded. Many species have broad bands of colour, the patterns helping to distinguish different species. Whereas honeybees have short tongues and therefore mainly pollinate open flowers, some bumblebee species have long tongues and collect nectar from flowers that are closed into a tube. Bumblebees have fewer stripes (or none), and usually have part of the body covered in black fur, while honeybees have many stripes including several grey stripes on the abdomen. Sizes are very variable even within species; the largest British species, B. terrestris, has queens up to long, males up to long, and workers between long. The largest bumblebee species in the world is B. dahlbomii of Chile, up to about long, and described as "flying mice" and "a monstrous fluffy ginger beast".
Distribution and habitat
Bumblebees are typically found in temperate climates, and are often found at higher latitudes and altitudes than other bees, although a few lowland tropical species exist. A few species (B. polaris and B. alpinus) range into very cold climates where other bees might not be found; B. polaris occurs in northern Ellesmere Island in the high Arctic, along with another bumblebee B. hyperboreus, which parasitises its nest. This is the northernmost occurrence of any eusocial insect. One reason for their presence in cold places is that bumblebees can regulate their body temperature, via solar radiation, internal mechanisms of "shivering" (called heterothermy), and countercurrent exchange to retain heat. Other bees have similar physiology, but the mechanisms seem best developed and have been most studied in bumblebees. They adapt to higher elevations by extending their wing stroke amplitude. Bumblebees have a largely cosmopolitan distribution but are absent from Australia (apart from Tasmania where they have been introduced) and are found in Africa only north of the Sahara. More than a hundred years ago they were also introduced to New Zealand, where they play an important role as efficient pollinators.
Biology
Feeding
The bumblebee tongue (the proboscis) is a long, hairy structure that extends from a sheath-like modified maxilla. The primary action of the tongue is lapping, that is, repeated dipping of the tongue into liquid.
The tip of the tongue probably acts as a suction cup and during lapping, nectar may be drawn up the proboscis by capillary action. When at rest or flying, the proboscis is kept folded under the head. The longer the tongue, the deeper the bumblebee can probe into a flower and bees probably learn from experience which flower source is best-suited to their tongue length. Bees with shorter proboscides, like Bombus bifarius, have a more difficult time foraging nectar relative to other bumblebees with longer proboscides; to overcome this disadvantage, B. bifarius workers were observed to lick the back of spurs on the nectar duct, which resulted in a small reward.
Wax production
The exoskeleton of the abdomen is divided into plates called dorsal tergites and ventral sternites. Wax is secreted from glands on the abdomen and extruded between the sternites where it resembles flakes of dandruff. It is secreted by the queen when she starts a nest and by young workers. It is scraped from the abdomen by the legs, moulded until malleable and used in the construction of honeypots, to cover the eggs, to line empty cocoons for use as storage containers and sometimes to cover the exterior of the nest.
Coloration
The brightly coloured pile of the bumblebee is an aposematic (warning) signal, given that females can inflict a painful sting. Depending on the species and morph, the warning colours range from entirely black, to bright yellow, red, orange, white, and pink. Dipteran flies in the families Syrphidae (hoverflies), Asilidae (robber flies), Tabanidae (horseflies), Oestridae (bot or warble flies) and Bombyliidae (bee flies, such as Bombylius major) all include Batesian mimics of bumblebees, resembling them closely enough to deceive at least some predators.
Many species of Bombus, including the group sometimes called Psithyrus (cuckoo bumblebees), have evolved Müllerian mimicry, where the different bumblebees in a region resemble each other, so that a young predator need only learn to avoid any of them once. For example, in California a group of bumblebees consists of largely black species including B. californicus, B. caliginosus, B. vandykei, B. vosnesenskii, B. insularis and B. fernaldae. Other bees in California include a group of species all banded black and yellow. In each case, Müllerian mimicry provides the bees in the group with a selective advantage. In addition, parasitic (cuckoo) bumblebees resemble their hosts more closely than would be expected by chance, at least in areas like Europe where parasite-host co-speciation is common; but this too may be explained as Müllerian mimicry, rather than requiring the parasite's coloration to deceive the host (aggressive mimicry).
Temperature control
Bumblebees are active under conditions during which honeybees and other smaller bees stay at home, and can readily absorb heat from even weak sunshine. The thick pile created by long setae (bristles) acts as insulation to keep bumblebees warm in cold weather; species from cold climates have longer setae (and thus thicker insulation) than those from the tropics. The temperature of the flight muscles, which occupy much of the thorax, needs to be at least before flight can take place. The muscle temperature can be raised by shivering. It takes about five minutes for the muscles to reach this temperature at an air temperature of .
Chill-coma temperature
The chill-coma temperature in relation to flying insects is the temperature at which flight muscles cannot be activated. Compared to honey bees and carpenter bees, bumblebees have the lowest chill-coma temperature. Of the bumblebees Bombus bimaculatus has the lowest at . However, bumblebees have been seen to fly in colder ambient temperatures. This discrepancy is likely because the chill-coma temperature was determined by tests done in a laboratory setting. However, bumblebees live in insulated shelters and can shiver to warm up before venturing into the cold.
Communication and social learning
Bumblebees do not have ears, and it is not known whether or how well they can hear. However, they are sensitive to the vibrations made by sound travelling through wood or other materials.
Bumblebees do not exhibit the "bee dances" used by honeybees to tell other workers the locations of food sources. Instead, when they return from a successful foraging expedition, they run excitedly around in the nest for several minutes before going out to forage once more. These bees may be offering some form of communication based on the buzzing sounds made by their wings, which may stimulate other bees to start foraging. Another stimulant to foraging activity is the level of food reserves in the colony. Bees monitor the amount of honey in the honeypots, and when little is left or when high-quality food is added, they are more likely to go out to forage.
Bumblebees have been observed to partake in social learning. In a 2017 study involving Bombus terrestris, bees were taught to complete an unnatural task of moving large objects to obtain a reward. Bees who first observed another bee complete the task were significantly more successful in learning the task than bees who observed the same action performed by a magnet, indicating the importance of social information. The bees did not copy one another exactly: in fact, the study suggested that the bees were instead attempting to emulate one another's goals.
Reproduction and nesting
Nest size depends on species of bumblebee. Most form colonies of between 50 and 400 individuals, but colonies have been documented as small as ~20 individuals and as large as 1700. These nests are small compared to honeybee hives, which hold about 50,000 bees. Many species nest underground, choosing old rodent burrows or sheltered places, and avoiding places that receive direct sunlight that could result in overheating. Other species make nests above ground, whether in thick grass or in holes in trees. A bumblebee nest is not organised into hexagonal combs like that of a honeybee; the cells are instead clustered together untidily. The workers remove dead bees or larvae from the nest and deposit them outside the nest entrance, helping to prevent disease. Nests in temperate regions last only for a single season and do not survive the winter.
In the early spring, the queen comes out of diapause and finds a suitable place to create her colony. Then she builds wax cells in which to lay her eggs which were fertilised the previous year. The eggs that hatch develop into female workers, and in time, the queen populates the colony, with workers feeding the young and performing other duties similar to honeybee workers. In temperate zones, young queens (gynes) leave the nest in the autumn and mate, often more than once, with males (drones) that are forcibly driven out of the colony. The drones and workers die as the weather turns colder; the young queens feed intensively to build up stores of fat for the winter. They survive in a resting state (diapause), generally below ground, until the weather warms up in the spring with the early bumblebee being the species that is among the first to emerge. Many species of bumblebee follow this general trend within the year. Bombus pensylvanicus is a species that follows this type of colony cycle. For this species the cycle begins in February, reproduction starts in July or August, and ends in the winter months. The queen remains in hibernation until spring of the following year in order to optimize conditions to search for a nest.
In fertilised queens, the ovaries only become active when the queen starts to lay. An egg passes along the oviduct to the vagina where there is a chamber called the spermatheca, in which the sperm from the mating is stored. Depending on need, she may allow her egg to be fertilised. Unfertilised eggs become haploid males; fertilised eggs grow into diploid females and queens. The hormones that stimulate the development of the ovaries are suppressed in female worker bees, while the queen remains dominant.
To develop, the larvae must be fed both nectar for carbohydrates and pollen for protein. Bumblebees feed nectar to the larvae by chewing a small hole in the brood cell into which they regurgitate nectar. Larvae are fed pollen in one of two ways, depending on the bumblebee species. Pocket-making bumblebees create pockets of pollen at the base of the brood-cell clump from which the larvae feed themselves. Pollen-storing bumblebees keep pollen in separate wax pots and feed it to the larvae.
After the emergence of the first or second group of offspring, workers take over the task of foraging and the queen spends most of her time laying eggs and caring for larvae. The colony grows progressively larger and eventually begins to produce males and new queens.
Bumblebee workers can lay unfertilised haploid eggs (with only a single set of chromosomes) that develop into viable male bumblebees. Only fertilised queens can lay diploid eggs (one set of chromosomes from a drone, one from the queen) that mature into workers and new queens.
In a young colony, the queen minimises reproductive competition from workers by suppressing their egg-laying through physical aggression and pheromones. Worker policing leads to nearly all eggs laid by workers being eaten. Thus, the queen is usually the mother of all of the first males laid. Workers eventually begin to lay male eggs later in the season when the queen's ability to suppress their reproduction diminishes. Because of the reproductive competition between workers and the queen, bumblebees are considered "primitively eusocial".
Although a large majority of bumblebees follow such monogynous colony cycles that only involve one queen, some select Bombus species (such as Bombus pauloensis) will spend part of their life cycle in a polygynous phase (have multiple queens in one nest during these periods of polygyny).
Foraging behaviour
Bumblebees generally visit flowers that exhibit the bee pollination syndrome and these patches of flowers may be up to 1–2 km from their colony. They tend to visit the same patches of flowers every day, as long as they continue to find nectar and pollen there, a habit known as pollinator or flower constancy. While foraging, bumblebees can reach ground speeds of up to .
Bumblebees use a combination of colour and spatial relationships to learn which flowers to forage from. They can also detect both the presence and the pattern of electric fields on flowers, which occur due to atmospheric electricity, and take a while to leak away into the ground. They use this information to find out if a flower has been recently visited by another bee. Bumblebees can detect the temperature of flowers, as well as which parts of the flower are hotter or cooler and use this information to recognise flowers. After arriving at a flower, they extract nectar using their long tongues ("glossae") and store it in their crops. Many species of bumblebees also exhibit "nectar robbing": instead of inserting the mouthparts into the flower in the normal way, these bees bite directly through the base of the corolla to extract nectar, avoiding pollen transfer.
Pollen is removed from flowers deliberately or incidentally by bumblebees. Incidental removal occurs when bumblebees come in contact with the anthers of a flower while collecting nectar. When it enters a flower, the bumblebee's body hairs receive a dusting of pollen from the anthers. In queens and workers this is then groomed into the corbiculae (pollen baskets) on the hind legs where it can be seen as bulging masses that may contain as many as a million pollen grains. Male bumblebees do not have corbiculae and do not purposively collect pollen. Bumblebees are also capable of buzz pollination, in which they dislodge pollen from the anthers by creating a resonant vibration with their flight muscles.
In at least some species, once a bumblebee has visited a flower, it leaves a scent mark on it. This scent mark deters bumblebees from visiting that flower until the scent degrades. This scent mark is a general chemical bouquet that bumblebees leave behind in different locations (e.g. nest, neutral, and food sites), and they learn to use this bouquet to identify both rewarding and unrewarding flowers, and may be able to identify who else has visited a flower. Bumblebees rely on this chemical bouquet more when the flower has a high handling time, that is, where it takes a longer time for the bee to find the nectar once inside the flower.
Once they have collected nectar and pollen, female workers return to the nest and deposit the harvest into brood cells, or into wax cells for storage. Unlike honeybees, bumblebees only store a few days' worth of food, so are much more vulnerable to food shortages. Male bumblebees collect only nectar and do so to feed themselves. They may visit quite different flowers from the workers because of their different nutritional needs.
Asynchronous flight muscles
Bees beat their wings about 200 times a second. Their thorax muscles do not contract on each nerve firing, but rather vibrate like a plucked rubber band. This is efficient, since it lets the system consisting of muscle and wing operate at its resonant frequency, leading to low energy consumption. Further, it is necessary, since insect motor nerves generally cannot fire 200 times per second. These types of muscles are called asynchronous muscles and are found in the insect wing systems in families such as Hymenoptera, Diptera, Coleoptera, and Hemiptera. Bumblebees must warm up their bodies considerably to get airborne at low ambient temperatures. Bumblebees can reach an internal thoracic temperature of 30 °C (86 °F) using this method.
Cuckoo bumblebees
Bumblebees of the subgenus Psithyrus (known as 'cuckoo bumblebees', and formerly considered a separate genus) are brood parasites, sometimes called kleptoparasites, in the colonies of other bumblebees, and have lost the ability to collect pollen. Before finding and invading a host colony, a Psithyrus female, such as that of the Psithyrus species of B. sylvestris, feeds directly from flowers. Once she has infiltrated a host colony, the Psithyrus female kills or subdues the queen of that colony, and uses pheromones and physical attacks to force the workers of that colony to feed her and her young. Usually, cuckoo bumblebees can be described as queen-intolerant inquilines, since the host queen is often killed to enable the parasite to produce more offspring, though some species, such as B. bohemicus, actually enjoy increased success when they leave the host queen alive.
The female Psithyrus has a number of morphological adaptations for combat, such as larger mandibles, a tough cuticle and a larger venom sac that increase her chances of taking over a nest. Upon emerging from their cocoons, the Psithyrus males and females disperse and mate. The males do not survive the winter but, like nonparasitic bumblebee queens, Psithyrus females find suitable locations to spend the winter and enter diapause after mating. They usually emerge from hibernation later than their host species. Each species of cuckoo bumblebee has a specific host species, which it may physically resemble. In the case of the parasitism of B. terrestris by B. (Psithyrus) vestalis, genetic analysis of individuals captured in the wild showed that about 42% of the host species' nests at a single location had "[lost] their fight against their parasite".
Sting
Queen and worker bumblebees can sting. Unlike in honeybees, a bumblebee's stinger lacks barbs, so the bee can sting repeatedly without leaving the stinger in the wound and thereby injuring itself. Bumblebee species are not normally aggressive, but may sting in defence of their nest, or if harmed. Female cuckoo bumblebees aggressively attack host colony members, and sting the host queen, but ignore other animals unless disturbed.
The sting is painful to humans, but not medically significant in most cases, although it may trigger an allergic reaction in susceptible individuals.
Predators, parasites, and pathogens
Bumblebees, despite their ability to sting, are eaten by certain predators. Nests may be dug up by badgers and eaten whole, including any adults present. Adults are preyed upon by robber flies and beewolves in North America. In Europe, birds including bee-eaters and shrikes capture adult bumblebees on the wing; smaller birds such as great tits also occasionally learn to take bumblebees, while camouflaged crab spiders catch them as they visit flowers.
The great grey shrike is able to detect flying bumblebees up to away; once captured, the sting is removed by repeatedly squeezing the insect with the mandibles and wiping the abdomen on a branch. The European honey buzzard follows flying bees back to their nest, digs out the nest with its feet, and eats larvae, pupae and adults as it finds them.
Bumblebees are parasitised by tracheal mites, Locustacarus buchneri; protozoans including Crithidia bombi and Apicystis bombi; and microsporidians including Nosema bombi and Nosema ceranae. The tree bumblebee B. hypnorum has spread into the United Kingdom despite hosting high levels of a nematode that normally interferes with queen bees' attempts to establish colonies. Deformed wing virus has been found to affect 11% of bumblebees in Great Britain.
Female bee moths (Aphomia sociella) prefer to lay their eggs in bumblebee nests. The A. sociella larvae will then feed on the eggs, larvae, and pupae left unprotected by the bumblebees, sometimes destroying large parts of the nest.
Relationship to humans
Agricultural use
Bumblebees are important pollinators of both crops and wildflowers. Because bumblebees do not overwinter the entire colony, they do not stockpile honey, and therefore are not useful as honey producers. Bumblebees are increasingly cultured for agricultural use as pollinators, among other reasons because they can pollinate plants such as tomato in greenhouses by buzz pollination whereas other pollinators cannot. Commercial production began in 1987, when Roland De Jonghe founded the Biobest company; in 1988 they produced enough nests to pollinate 40 hectares of tomatoes. The industry grew quickly, starting with other companies in the Netherlands. Bumblebee nests, mainly of buff-tailed bumblebees, are produced in at least 30 factories around the world; over a million nests are grown annually in Europe; Turkey is a major producer.
Bumblebees are Northern Hemisphere animals. When red clover was introduced as a crop to New Zealand in the nineteenth century, it was found to have no local pollinators, and clover seed had accordingly to be imported each year. Four species of bumblebee from the United Kingdom were therefore imported as pollinators. In 1885 and 1886, the Canterbury Acclimatization Society brought in 442 queens, of which 93 survived and quickly multiplied. As planned, red clover was soon being produced from locally-grown seed. Bumblebees are also reared commercially to pollinate tomatoes grown in greenhouses. The New Zealand population of buff-tailed bumblebees began colonising Tasmania, away, after being introduced there in 1992 under unclear circumstances.
Some concerns exist about the impact of the international trade in mass-produced bumblebee colonies. Evidence from Japan and South America indicates bumblebees can escape and naturalise in new environments, causing damage to native pollinators. Greater use of native pollinators, such as Bombus ignitus in China and Japan, has occurred as a result. In addition, mounting evidence indicates mass-produced bumblebees may also carry diseases, harmful to wild bumblebees and honeybees.
In Canada and Sweden, it has been shown that growing a mosaic of different crops encourages bumblebees and provides higher yields than does a monoculture of oilseed rape, despite the fact that the bees were attracted to the crop.
Population decline
Bumblebee species are declining in Europe, North America, and Asia due to a number of factors, including land-use change that reduces their food plants. In North America, pathogens are possibly having a stronger negative effect especially for the subgenus Bombus. A major impact on bumblebees was caused by the mechanisation of agriculture, accelerated by the urgent need to increase food production during the Second World War. Small farms depended on horses to pull implements and carts. The horses were fed on clover and hay, both of which were permanently grown on a typical farm. Little artificial fertiliser was used. Farms thus provided flowering clover and flower-rich meadows, favouring bumblebees. Mechanisation removed the need for horses and most of the clover; artificial fertilisers encouraged the growth of taller grasses, outcompeting the meadow flowers. Most of the flowers, and the bumblebees that fed on them, disappeared from Britain by the early 1980s. The last native British short-haired bumblebee was captured near Dungeness in 1988. This significant increase in pesticide and fertilizer use associated with the industrialization of agriculture has had adverse effects on the genus Bombus. The bees are directly exposed to the chemicals in two ways: by consuming nectar that has been directly treated with pesticide, or through physical contact with treated plants and flowers. The species Bombus hortorum in particular has been found to be affected by the pesticides; their brood development has been reduced and their memory has been negatively affected. Additionally, pesticide use negatively affects colony development and size.
Bumblebees are in danger in many developed countries due to habitat destruction and collateral pesticide damage. The European Food Safety Authority ruled that three neonicotinoid pesticides (clothianidin, imidacloprid, and thiamethoxam) presented a high risk for bees. While most work on neonicotinoid toxicity has looked at honeybees, a study on B. terrestris showed that "field-realistic" levels of imidacloprid significantly reduced growth rate and cut production of new queens by 85%, implying a "considerable negative effect" on wild bumblebee populations throughout the developed world. Another study on B. terrestris had results suggesting that use of neonicotinoid pesticides can affect how well bumblebees are able to forage and pollinate. Foragers from bee colonies that had been affected by the pesticide took longer to learn to manipulate flowers and visited flowers with less nutritious pollen.
In another study, chronic exposure in a laboratory setting to field-realistic levels of the neonicotinoid pesticide thiamethoxam did not affect colony weight gain or the number or mass of sexuals produced. Low levels of neonicotinoids can reduce the number of bumblebees in a colony by as much as 55%, and cause dysfunction in the bumblebees' brains. The Bumblebee Conservation Trust considers this evidence of reduced brain function "particularly alarming given that bumblebees rely upon their intelligence to go about their daily tasks."
Of 19 species of native nestmaking bumblebees and six species of cuckoo bumblebees formerly widespread in Britain, three have been extirpated, eight are in serious decline, and only six remain widespread. Similar declines have been reported in Ireland, with four species designated endangered, and another two considered vulnerable to extinction. A decline in bumblebee numbers could cause large-scale changes to the countryside, resulting from inadequate pollination of certain plants.
Some bumblebees native to North America are also vanishing, such as Bombus balteatus, Bombus terricola, Bombus affinis, and Bombus occidentalis; one, Bombus franklini, may be extinct. In South America, Bombus bellicosus was extirpated in the northern limit of its distribution range, probably due to intense land use and climate change effects.
Conservation efforts
In 2006, the bumblebee researcher Dave Goulson founded a registered charity, the Bumblebee Conservation Trust, to prevent the extinction "of any of the UK's bumblebees." In 2009 and 2010, the Trust attempted to reintroduce the short-haired bumblebee, Bombus subterraneus, which had become extinct in Britain, from the British-derived populations surviving in New Zealand from their introduction there a century earlier. From 2011, the Trust, in partnership with Natural England, Hymettus and the RSPB, has reintroduced short-haired bumblebee queens from Skåne in southern Sweden to restored flower-rich meadows at Dungeness in Kent. The queens were checked for mites and American foulbrood disease. Agri-environment schemes spread across the neighbouring area of Romney Marsh have been set up to provide over 800 hectares of additional flower-rich habitat for the bees. By the summer of 2013, workers of the species were found near the release zone, proving that nests had been established. The restored habitat has produced a revival in at least five "Schedule 41 priority" species: the ruderal bumblebee, Bombus ruderatus; the red-shanked carder bee, Bombus ruderarius; the shrill carder bee, Bombus sylvarum; the brown-banded carder bee, Bombus humilis and the moss carder bee, Bombus muscorum.
The world's first bumblebee sanctuary was established at Vane Farm in the Loch Leven National Nature Reserve in Scotland in 2008. In 2011, London's Natural History Museum led the establishment of an International Union for Conservation of Nature Bumblebee Specialist Group, chaired by Dr. Paul H. Williams, to assess the threat status of bumblebee species worldwide using Red List criteria.
Bumblebee conservation is in its infancy in many parts of the world, but with the realization of the important part they play in pollination of crops, efforts are being made to manage farmland better. Enhancing the wild bee population can be done by the planting of wildflower strips, and in New Zealand, bee nesting boxes have achieved some success, perhaps because there are few burrowing mammals to provide potential nesting sites in that country.
Misconception about flight
According to 20th-century folklore, the laws of aerodynamics prove the bumblebee should be incapable of flight, as it does not have the capacity (in terms of wing size or beats per second) to achieve flight with the degree of wing loading necessary.
The origin of this claim has been difficult to pin down with any certainty. John H. McMasters recounted an anecdote about an unnamed Swiss aerodynamicist at a dinner party who performed some rough calculations and concluded, presumably in jest, that according to the equations, bumblebees cannot fly.
In later years, McMasters backed away from this origin, suggesting there could be multiple sources, and the earliest he has found was a reference in the 1934 book by French entomologist Antoine Magnan (1881–1938); they had applied the equations of air resistance to insects and found their flight was impossible, but "One shouldn't be surprised that the results of the calculations don't square with reality". The following passage appears in the introduction to Le Vol des Insectes:
Magnan refers to his assistant André Sainte-Laguë.
Some credit physicist Ludwig Prandtl (1875–1953) of the University of Göttingen in Germany with popularizing the idea. Others say Swiss gas dynamicist Jakob Ackeret (1898–1981) did the calculations.
The calculations that purported to show that bumblebees cannot fly are based upon a simplified linear treatment of oscillating aerofoils. The method assumes small amplitude oscillations without flow separation. This ignores the effect of dynamic stall (an airflow separation inducing a large vortex above the wing), which briefly produces several times the lift of the aerofoil in regular flight. More sophisticated aerodynamic analysis shows the bumblebee can fly because its wings encounter dynamic stall in every oscillation cycle.
Additionally, John Maynard Smith, a noted biologist with a strong background in aeronautics, has pointed out that bumblebees would not be expected to sustain flight, as they would need to generate too much power given their tiny wing area. However, in aerodynamics experiments with other insects, he found that viscosity at the scale of small insects meant even their small wings can move a very large volume of air relative to their size, and this reduces the power required to sustain flight by an order of magnitude.
In music and literature
The orchestral interlude Flight of the Bumblebee was composed (c. 1900) by Nikolai Rimsky-Korsakov. It represents the turning of Prince Guidon into a bumblebee so he can fly away to visit his father, Tsar Saltan, in the opera The Tale of Tsar Saltan, although the music may reflect the flight of a bluebottle rather than a bumblebee. The music inspired Walt Disney to feature a bumblebee in his 1940 animated musical Fantasia and have it sound as if it were flying in all parts of the theater. This early attempt at "surround sound" was excluded from the film in later showings.
In 1599, during the reign of Queen Elizabeth I, someone, possibly Tailboys Dymoke, published Caltha Poetarum: Or The Bumble Bee, under the pseudonym "T. Cutwode". This was one of nine books censored under the Bishops' Ban issued by the Archbishop of Canterbury John Whitgift and the Bishop of London Richard Bancroft.
Emily Dickinson made a bumblebee the subject of her parody of Isaac Watts's well-known poem about honeybees, "How Doth the Little Busy Bee" (1715). Where Watts wrote "How skilfully she builds her cell! How neat she spreads the wax!", Dickinson's poem, "The Bumble-Bee's Religion" (1881), begins "His little Hearse-like Figure / Unto itself a Dirge / To a delusive Lilac / The vanity divulge / Of Industry and Morals / And every righteous thing / For the divine Perdition / of Idleness and Spring." The letter was said to have enclosed a dead bee.
In 1847, Ralph Waldo Emerson published his poem "The Humble-Bee".
The entomologist Otto Plath wrote Bumblebees and Their Ways in 1934. His daughter, the poet Sylvia Plath, wrote a group of poems about bees late in 1962, within four months of her suicide, transforming her father's interest into her poetry.
The scientist and illustrator Moses Harris (1731–1785) painted accurate watercolour drawings of bumblebees in his An Exposition of English Insects Including the Several Classes of Neuroptera, Hymenoptera, & Diptera, or Bees, Flies, & Libellulae (1776–80).
Bumblebees appear as characters, often eponymously, in children's books. The surname Dumbledore in the Harry Potter series (1997–2007) is an old name for bumblebee. J. K. Rowling said the name "seemed to suit the headmaster, because one of his passions is music and I imagined him walking around humming to himself". J. R. R. Tolkien, in his poem Errantry, also used the name Dumbledor, but for a large bee-like creature.
Among the many books for younger children are Bumble the Bee by Yvon Douran and Tony Neal (2014); Bertie Bumble Bee by K. I. Al-Ghani (2012); Ben the Bumble Bee: How do bees make honey? by Romessa Awadalla (2015); Bumble Bee Bob Has a Big Butt by Papa Campbell (2012); Buzz, Buzz, Buzz! Went Bumble-bee by Colin West (1997); Bumble Bee by Margaret Wise Brown (2000); How the Bumble Came to Bee by Paul and Ella Quarry (2012); The Adventures of Professor Bumble and the Bumble Bees by Stephen Brailovsky (2010). Among Beatrix Potter's "little books", Babbity Bumble and other members of her nest appear in The Tale of Mrs. Tittlemouse (1910).
Military
The United States Naval Construction Battalions adopted the bumblebee as their insignia in 1942.
| Biology and health sciences | Hymenoptera | null |
197179 | https://en.wikipedia.org/wiki/Sexual%20dimorphism | Sexual dimorphism | Sexual dimorphism is the condition where sexes of the same species exhibit different morphological characteristics, including characteristics not directly involved in reproduction. The condition occurs in most dioecious species, which consist of most animals and some plants. Differences may include secondary sex characteristics, size, weight, color, markings, or behavioral or cognitive traits. Male-male reproductive competition has evolved a diverse array of sexually dimorphic traits. Aggressive utility traits such as "battle" teeth and blunt heads reinforced as battering rams are used as weapons in aggressive interactions between rivals. Passive displays such as ornamental feathering or song-calling have also evolved mainly through sexual selection. These differences may be subtle or exaggerated and may be subjected to sexual selection and natural selection. The opposite of dimorphism is monomorphism, when both biological sexes are phenotypically indistinguishable from each other.
Overview
Ornamentation and coloration
Common and easily identified types of dimorphism consist of ornamentation and coloration, though not always apparent. A difference in the coloration of sexes within a given species is called sexual dichromatism, commonly seen in many species of birds and reptiles. Sexual selection leads to exaggerated dimorphic traits that are used predominantly in competition over mates. The increased fitness resulting from ornamentation offsets its cost to produce or maintain, suggesting complex evolutionary implications, but the costs and evolutionary implications vary from species to species.
The peafowl constitute conspicuous illustrations of the principle. The ornate plumage of peacocks, as used in the courting display, attracts peahens. At first sight, one might mistake peacocks and peahens for completely different species because of the vibrant colours and the sheer size of the male's plumage; the peahen is of a subdued brown coloration. The plumage of the peacock increases its vulnerability to predators because it is a hindrance in flight, and it renders the bird conspicuous in general. Similar examples are manifold, such as in birds of paradise and argus pheasants.
Another example of sexual dichromatism is that of nestling blue tits. Males are chromatically more yellow than females. It is believed that this is obtained by the ingestion of green Lepidopteran larvae, which contain large amounts of the carotenoids lutein and zeaxanthin. This diet also affects the sexually dimorphic colours in the human-invisible ultraviolet spectrum. Hence, the male birds, although appearing yellow to humans, actually have a violet-tinted plumage that is seen by females. This plumage is thought to be an indicator of male parental abilities. Perhaps this is a good indicator for females because it shows that they are good at obtaining a food supply from which the carotenoid is obtained. There is a positive correlation between the chromas of the tail and breast feathers and body condition. Carotenoids play an important role in immune function for many animals, so carotenoid dependent signals might indicate health.
Frogs constitute another conspicuous illustration of the principle. There are two types of dichromatism for frog species: ontogenetic and dynamic. Ontogenetic frogs are more common and have permanent color changes in males or females. Ranoidea lesueuri is an example of a dynamic frog with temporary color changes in males during the breeding season. Hyperolius ocellatus is an ontogenetic frog with dramatic differences in both color and pattern between the sexes. At sexual maturity, the males display a bright green with white dorsolateral lines. In contrast, the females are rusty red to silver with small spots. The bright coloration in the male population attracts females and is an aposematic sign to potential predators.
Females often show a preference for exaggerated male secondary sexual characteristics in mate selection. The sexy son hypothesis explains that females prefer more elaborate males and select against males that are dull in color, independent of the species' vision.
Similar sexual dimorphism and mating choice are also observed in many fish species. For example, male guppies have colorful spots and ornamentations, while females are generally grey. Female guppies prefer brightly colored males to duller males.
In redlip blennies, only the male fish develops an organ at the anal-urogenital region that produces antimicrobial substances. During parental care, males rub their anal-urogenital regions over their nests' internal surfaces, thereby protecting their eggs from microbial infections, one of the most common causes for mortality in young fish.
Plants
Most flowering plants are hermaphroditic but approximately 6% of species have separate males and females (dioecy). Sexual dimorphism is common in dioecious plants and dioicous species.
Males and females in insect-pollinated species generally look similar to one another because plants provide rewards (e.g. nectar) that encourage pollinators to visit another similar flower, completing pollination. Catasetum orchids are one interesting exception to this rule. Male Catasetum orchids violently attach pollinia to euglossine bee pollinators. The bees will then avoid other male flowers but may visit the female, which looks different from the males.
Various other dioecious exceptions, such as Loxostylis alata have visibly different sexes, with the effect of eliciting the most efficient behavior from pollinators, who then use the most efficient strategy in visiting each gender of flower instead of searching, say, for pollen in a nectar-bearing female flower.
Some plants, such as some species of Geranium have what amounts to serial sexual dimorphism. The flowers of such species might, for example, present their anthers on opening, then shed the exhausted anthers after a day or two and perhaps change their colours as well while the pistil matures; specialist pollinators are very much inclined to concentrate on the exact appearance of the flowers they serve, which saves their time and effort and serves the interests of the plant accordingly. Some such plants go even further and change their appearance once fertilized, thereby discouraging further visits from pollinators. This is advantageous to both parties because it avoids damaging the developing fruit and wasting the pollinator's effort on unrewarding visits. In effect, the strategy ensures that pollinators can expect a reward every time they visit an appropriately advertising flower.
Females of the aquatic plant Vallisneria americana have floating flowers attached by a long flower stalk that are fertilized if they contact one of the thousands of free-floating flowers released by a male. Sexual dimorphism is most often associated with wind-pollination in plants due to selection for efficient pollen dispersal in males vs pollen capture in females, e.g. Leucadendron rubrum.
Sexual dimorphism in plants can also be dependent on reproductive development. This can be seen in Cannabis sativa, a type of hemp, which have higher photosynthesis rates in males while growing but higher rates in females once the plants become sexually mature.
Every sexually reproducing extant species of the vascular plant has an alternation of generations; the plants we see about us generally are diploid sporophytes, but their offspring are not the seeds that people commonly recognize as the new generation. The seed actually is the offspring of the haploid generation of microgametophytes (pollen) and megagametophytes (the embryo sacs in the ovules). Each pollen grain accordingly may be seen as a male plant in its own right; it produces a sperm cell and is dramatically different from the female plant, the megagametophyte that produces the female gamete.
Insects
Insects display a wide variety of sexual dimorphism between taxa including size, ornamentation and coloration. The female-biased sexual size dimorphism observed in many taxa evolved despite intense male-male competition for mates. In Osmia rufa, for example, the female is larger/broader than males, with males being 8–10 mm in size and females being 10–12 mm in size. In the hackberry emperor females are similarly larger than males. The reason for the sexual dimorphism is due to provision size mass, in which females consume more pollen than males.
In some species, there is evidence of male dimorphism, but it appears to be for distinctions of roles. This is seen in the bee species Macrotera portalis in which there is a small-headed morph, capable of flight, and large-headed morph, incapable of flight, for males. Anthidium manicatum also displays male-biased sexual dimorphism. The selection for larger size in males rather than females in this species may have resulted due to their aggressive territorial behavior and subsequent differential mating success. Another example is Lasioglossum hemichalceum, which is a species of sweat bee that shows drastic physical dimorphisms between male offspring. Not all dimorphism has to have a drastic difference between the sexes. Andrena agilissima is a mining bee where the females only have a slightly larger head than the males.
Weaponry leads to increased fitness by increasing success in male–male competition in many insect species. The beetle horns in Onthophagus taurus are enlarged growths of the head or thorax expressed only in the males. Copris ochus also has distinct sexual and male dimorphism in head horns. Another beetle with a distinct horn-related sexual dimorphism is Allomyrina dichotoma, also known as the Japanese rhinoceros beetle. These structures are impressive because of the exaggerated sizes. There is a direct correlation between male horn lengths and body size and higher access to mates and fitness. In other beetle species, both males and females may have ornamentation such as horns.
Generally, insect sexual size dimorphism (SSD) within species increases with body size.
Sexual dimorphism within insects is also displayed by dichromatism. In butterfly genera Bicyclus and Junonia, dimorphic wing patterns evolved due to sex-limited expression, which mediates the intralocus sexual conflict and leads to increased fitness in males. The sexual dichromatic nature of Bicyclus anynana is reflected by female selection on the basis of dorsal UV-reflective eyespot pupils. The common brimstone also displays sexual dichromatism; males have yellow and iridescent wings, while female wings are white and non-iridescent. Naturally selected deviation in protective female coloration is displayed in mimetic butterflies.
Spiders and sexual cannibalism
Many arachnid groups exhibit sexual dimorphism, but it is most widely studied in the spiders. In the orb-weaving spider Zygiella x-notata, for example, adult females have a larger body size than adult males. Size dimorphism shows a correlation with sexual cannibalism, which is prominent in spiders (it is also found in insects such as praying mantises). In the size dimorphic wolf spider Tigrosa helluo, food-limited females cannibalize more frequently. Therefore, there is a high risk of low fitness for males due to pre-copulatory cannibalism, which led to male selection of larger females for two reasons: higher fecundity and lower rates of cannibalism. In addition, female fecundity is positively correlated with female body size and large female body size is selected for, which is seen in the family Araneidae. All Argiope species, including Argiope bruennichi, use this method. Some males evolved ornamentation including binding the female with silk, having proportionally longer legs, modifying the female's web, mating while the female is feeding, or providing a nuptial gift in response to sexual cannibalism. Male body size is not under selection due to cannibalism in all spider species such as Nephila pilipes, but is more prominently selected for in less dimorphic species of spiders, which often selects for larger male size. In the species Maratus volans, the males are known for their characteristic colorful fan which attracts the females during mating.
Fish
Ray-finned fish are an ancient and diverse class, with the widest degree of sexual dimorphism of any animal class. Fairbairn notes that "females are generally larger than males but males are often larger in species with male–male combat or male paternal care ... [sizes range] from dwarf males to males more than 12 times heavier than females."
There are cases where males are substantially larger than females. An example is Lamprologus callipterus, a type of cichlid fish. In this fish, the males are characterized as being up to 60 times larger than the females. The male's increased size is believed to be advantageous because males collect and defend empty snail shells in each of which a female breeds. Males must be larger and more powerful in order to collect the largest shells. The female's body size must remain small because in order for her to breed, she must lay her eggs inside the empty shells. If she grows too large, she will not fit in the shells and will be unable to breed. The female's small body size is also likely beneficial to her chances of finding an unoccupied shell. Larger shells, although preferred by females, are often limited in availability. Hence, the female is limited to the growth of the size of the shell and may actually change her growth rate according to shell size availability. In other words, the male's ability to collect large shells depends on his size. The larger the male, the larger the shells he is able to collect. This then allows for females to be larger in his brooding nest which makes the difference between the sizes of the sexes less substantial. Male–male competition in this fish species also selects for large size in males. There is aggressive competition by males over territory and access to larger shells. Large males win fights and steal shells from competitors. Another example is the dragonet, in which males are considerably larger than females and possess longer fins.
Sexual dimorphism also occurs in hermaphroditic fish. These species are known as sequential hermaphrodites. In fish, reproductive histories often include the sex-change from female to male where there is a strong connection between growth, the sex of an individual, and the mating system within which it operates. In protogynous mating systems where males dominate mating with many females, size plays a significant role in male reproductive success. Males have a propensity to be larger than females of a comparable age but it is unclear whether the size increase is due to a growth spurt at the time of the sexual transition or due to the history of faster growth in sex changing individuals. Larger males are able to stifle the growth of females and control environmental resources.
Social organization plays a large role in the changing of sex by the fish. It is often seen that a fish will change its sex when there is a lack of a dominant male within the social hierarchy. The females that change sex are often those who attain and preserve an initial size advantage early in life. In either case, females which change sex to males are larger and often prove to be a good example of dimorphism.
In other cases with fish, males will go through noticeable changes in body size, and females will go through morphological changes that can only be seen inside of the body. For example, in sockeye salmon, males develop larger body size at maturity, including an increase in body depth, hump height, and snout length. Females experience minor changes in snout length, but the most noticeable difference is the huge increase in gonad size, which accounts for about 25% of body mass.
Sexual selection was observed for female ornamentation in Gobiusculus flavescens, known as two-spotted gobies. Traditional hypotheses suggest that male–male competition drives selection. However, selection for ornamentation within this species suggests that showy female traits can be selected through either female–female competition or male mate choice. Since carotenoid-based ornamentation suggests mate quality, female two-spotted guppies that develop colorful orange bellies during breeding season are considered favorable to males. The males invest heavily in offspring during incubation, which leads to the sexual preference in colorful females due to higher egg quality.
Amphibians and non-avian reptiles
In amphibians and reptiles, the degree of sexual dimorphism varies widely among taxonomic groups. The sexual dimorphism in amphibians and reptiles may be reflected in any of the following: anatomy; relative length of tail; relative size of head; overall size as in many species of vipers and lizards; coloration as in many amphibians, snakes, and lizards, as well as in some turtles; an ornament as in many newts and lizards; the presence of specific sex-related behaviour is common to many lizards; and vocal qualities which are frequently observed in frogs.
Anole lizards show prominent size dimorphism with males typically being significantly larger than females. For instance, the average male Anolis sagrei was 53.4 mm vs. 40 mm in females. Different sizes of the heads in anoles have been explained by differences in the estrogen pathway. The sexual dimorphism in lizards is generally attributed to the effects of sexual selection, but other mechanisms including ecological divergence and fecundity selection provide alternative explanations. The development of color dimorphism in lizards is induced by hormonal changes at the onset of sexual maturity, as seen in Psamodromus algirus, Sceloporus gadoviae, and S. undulates erythrocheilus. Sexual dimorphism in size is also seen in frog species like P. bibronii.
Male painted dragon lizards, Ctenophorus pictus. are brightly conspicuous in their breeding coloration, but male colour declines with aging. Male coloration appears to reflect innate anti-oxidation capacity that protects against oxidative DNA damage. Male breeding coloration is likely an indicator to females of the underlying level of oxidative DNA damage (a significant component of aging) in potential mates.
Birds
Possible mechanisms have been proposed to explain macroevolution of sexual size dimorphism in birds. These include sexual selection, selection for fecundity in females, niche divergence between the sexes, and allometry, but their relative importance is still not fully understood . Sexual dimorphism in birds can be manifested in size or plumage differences between the sexes. Sexual size dimorphism varies among taxa, with males typically being larger, though this is not always the case, e.g. birds of prey, hummingbirds, and some species of flightless birds. Plumage dimorphism, in the form of ornamentation or coloration, also varies, though males are typically the more ornamented or brightly colored sex. Such differences have been attributed to the unequal reproductive contributions of the sexes. This difference produces a stronger female choice since they have more risk in producing offspring. In some species, the male's contribution to reproduction ends at copulation, while in other species the male becomes the main (or only) caregiver. Plumage polymorphisms have evolved to reflect these differences and other measures of reproductive fitness, such as body condition or survival. The male phenotype sends signals to females who then choose the 'fittest' available male.
Sexual dimorphism is a product of both genetics and environmental factors. An example of sexual polymorphism determined by environmental conditions exists in the red-backed fairywren. Red-backed fairywren males can be classified into three categories during breeding season: black breeders, brown breeders, and brown auxiliaries. These differences arise in response to the bird's body condition: if they are healthy they will produce more androgens thus becoming black breeders, while less healthy birds produce less androgens and become brown auxiliaries. The reproductive success of the male is thus determined by his success during each year's non-breeding season, causing reproductive success to vary with each year's environmental conditions.
Migratory patterns and behaviors also influence sexual dimorphisms. This aspect also stems back to size dimorphism in species. It has been shown that the larger males are better at coping with the difficulties of migration and thus are more successful in reproducing when reaching the breeding destination. When viewing this from an evolutionary standpoint, many theories and explanations come into consideration. If these are the result for every migration and breeding season, the expected results should be a shift towards a larger male population through sexual selection. Sexual selection is strong when the factor of environmental selection is also introduced. Environmental selection may support a smaller chick size if those chicks were born in an area that allowed them to grow to a larger size, even though under normal conditions they would not be able to reach this optimal size for migration. When the environment gives advantages and disadvantages of this sort, the strength of selection is weakened and the environmental forces are given greater morphological weight. The sexual dimorphism could also produce a change in timing of migration leading to differences in mating success within the bird population. When the dimorphism produces that large of a variation between the sexes and between the members of the sexes, multiple evolutionary effects can take place. This timing could even lead to a speciation phenomenon if the variation becomes strongly drastic and favorable towards two different outcomes. Sexual dimorphism is maintained by the counteracting pressures of natural selection and sexual selection. For example, sexual dimorphism in coloration increases the vulnerability of bird species to predation by European sparrowhawks in Denmark. Presumably, increased sexual dimorphism means males are brighter and more conspicuous, leading to increased predation. Moreover, the production of more exaggerated ornaments in males may come at the cost of suppressed immune function. So long as the reproductive benefits of the trait due to sexual selection are greater than the costs imposed by natural selection, then the trait will propagate throughout the population. Reproductive benefits arise in the form of a larger number of offspring, while natural selection imposes costs in the form of reduced survival. This means that even if the trait causes males to die earlier, the trait is still beneficial so long as males with the trait produce more offspring than males lacking the trait. This balance keeps dimorphism alive in these species and ensures that the next generation of successful males will also display these traits that are attractive to females.
Such differences in form and reproductive roles often cause differences in behavior. As previously stated, males and females often have different roles in reproduction. The courtship and mating behavior of males and females are regulated largely by hormones throughout a bird's lifetime. Activational hormones occur during puberty and adulthood and serve to 'activate' certain behaviors when appropriate, such as territoriality during breeding season. Organizational hormones occur only during a critical period early in development, either just before or just after hatching in most birds, and determine patterns of behavior for the rest of the bird's life. Such behavioral differences can cause disproportionate sensitivities to anthropogenic pressures. Females of the whinchat in Switzerland breed in intensely managed grasslands. Earlier harvesting of the grasses during the breeding season lead to more female deaths. Populations of many birds are often male-skewed and when sexual differences in behavior increase this ratio, populations decline at a more rapid rate. Also not all male dimorphic traits are due to hormones like testosterone, instead they are a naturally occurring part of development, for example plumage. In addition, the strong hormonal influence on phenotypic differences suggests that the genetic mechanism and genetic basis of these sexually dimorphic traits may involve transcription factors or cofactors rather than regulatory sequences.
Sexual dimorphism may also influence differences in parental investment during times of food scarcity. For example, in the blue-footed booby, the female chicks grow faster than the males, resulting in booby parents producing the smaller sex, the males, during times of food shortage. This then results in the maximization of parental lifetime reproductive success. In Black-tailed Godwits Limosa limosa limosa females are also the larger sex, and the growth rates of female chicks are more susceptible to limited environmental conditions.
Sexual dimorphism may also only appear during mating season; some species of birds only show dimorphic traits in seasonal variation. The males of these species will molt into a less bright or less exaggerated color during the off-breeding season. This occurs because the species is more focused on survival than on reproduction, causing a shift into a less ornate state.
Consequently, sexual dimorphism has important ramifications for conservation. However, sexual dimorphism is not only found in birds and is thus important to the conservation of many animals. Such differences in form and behavior can lead to sexual segregation, defined as sex differences in space and resource use. Most sexual segregation research has been done on ungulates, but such research extends to bats, kangaroos, and birds. Sex-specific conservation plans have even been suggested for species with pronounced sexual segregation.
The term sesquimorphism (the Latin numeral prefix sesqui- means one-and-one-half, so halfway between mono- (one) and di- (two)) has been proposed for bird species in which "both sexes have basically the same plumage pattern, though the female is clearly distinguishable by reason of her paler or washed-out Examples include Cape sparrow (Passer melanurus), rufous sparrow (subspecies P. motinensis motinensis), and saxaul sparrow (P. ammodendri).
Non-avian dinosaurs
Examining fossils of non-avian dinosaurs in search of sexually dimorphic characteristics requires the supply of complete and articulated skeletal and tissue remains. As terrestrial organisms, dinosaur carcasses are subject to ecological and geographical influence that inevitably constitutes the degree of preservation. The availability of well-preserved remains is not a probable outcome as a consequence of decomposition and fossilization. Some paleontologists have looked for sexual dimorphism among dinosaurs using statistics and comparison to ecologically or phylogenetically related modern animals.
Apatosaurus and Diplodocus
Female Apatosaurus and Diplodocus had interconnected caudal vertebrae that allowed them to keep their tails elevated to aid in copulation. Discovering that this fusion occurred in only 50% of Apatosaurus and Diplodocus skeletons and 25% of Camarasaurus skeletons indicated that this is a sexually dimorphic trait.
Theropoda
It has been hypothesized that male theropods possessed a retractable penis, a feature similar to modern day crocodilians. Crocodilian skeletons were examined to determine whether there is a skeletal component that is distinctive between both sexes, to help provide an insight on the physical disparities between male and female theropods. Findings revealed the caudal chevrons of male crocodiles, used to anchor the penis muscles, were significantly larger than those of females. There have been criticisms of these findings, but it remains a subject of debate among advocates and adversaries.
Ornithopoda
Studies of sexual dimorphism in hadrosaurs have generally centered on the distinctive cranial crests, which likely provided a function in sexual display. A biometric study of 36 skulls found sexual dimorphism was exhibited in the crest of 3 species of hadrosaurids. The crests could be categorized as full (male) or narrow (female) and may have given some advantage in intrasexual mating-competition.
Ceratopsians
According to Scott D. Sampson, if ceratopsids were to exhibit sexual dimorphism, modern ecological analogues suggest it would be found in display structures, such as horns and frills. No convincing evidence for sexual dimorphism in body size or mating signals is known in ceratopsids, although there is evidence that the more primitive ceratopsian Protoceratops andrewsi possessed sexes that were distinguishable based on frill and nasal prominence size. This is consistent with other known tetrapod groups where midsized animals tend to exhibit markedly more sexual dimorphism than larger ones. However, it has been proposed that these differences can be better explained by intraspecific and ontogenic variation rather than sexual dimorphism. In addition, many sexually dimorphic traits that may have existed in ceratopsians include soft tissue variations such as coloration or dewlaps, which would be unlikely to have been preserved in the fossil record.
Stegosaurians
A 2015 study on specimens of Hesperosaurus mjosi found evidence of sexual dimorphism in the shape of the dermal plates. Two plate morphs were described: one was short, wide, and oval-shaped, the other taller and narrower.
Mammals
In 45% of mammal species, males are larger than females, in 39%, males and females are the same size and in 16% of mammals, females are larger than males. Both genes and hormones affect the formation of many animal brains before "birth" (or hatching), and also behaviour of adult individuals. Hormones significantly affect human brain formation, and also brain development at puberty. A 2004 review in Nature Reviews Neuroscience observed that "because it is easier to manipulate hormone levels than the expression of sex chromosome genes, the effects of hormones have been studied much more extensively, and are much better understood, than the direct actions in the brain of sex chromosome genes." It concluded that while "the differentiating effects of gonadal secretions seem to be dominant," the existing body of research "support the idea that sex differences in neural expression of X and Y genes significantly contribute to sex differences in brain functions and disease."
Pinnipeds
Marine mammals show some of the greatest sexual size differences of mammals, because of sexual selection and environmental factors like breeding location. The mating system of pinnipeds varies from polygamy to serial monogamy. Pinnipeds are known for early differential growth and maternal investment since the only nutrients for newborn pups is the milk provided by the mother. For example, the males are significantly larger (about 10% heavier and 2% longer) than the females at birth in sea lion pups. The pattern of differential investment can be varied principally prenatally and post-natally. Mirounga leonina, the southern elephant seal, is one of the most dimorphic mammals.
Primates
Humans
According to Clark Spencer Larsen, modern day Homo sapiens show a range of sexual dimorphism, with average body mass between the sexes differing by roughly 15%. Substantial discussion in academic literature considers potential evolutionary advantages associated with sexual competition (both intrasexual and intersexual), as well as short- and long-term sexual strategies. According to Daly and Wilson, "The sexes differ more in human beings than in monogamous mammals, but much less than in extremely polygamous mammals."
The average basal metabolic rate is about 6 percent higher in adolescent males than females and increases to about 10 percent higher after puberty. Females tend to convert more food into fat, while males convert more into muscle and expendable circulating energy reserves. According to Tim Hewett, director of research in the department of sports medicine at Ohio State University Wexner Medical Center, females have, on average, 50–60% of the upper body strength of males, and 80–90% of the lower body strength of males, relative to body size.
The difference in strength relative to body mass is less pronounced in trained individuals. In Olympic weightlifting, male records vary from 5.5× body mass in the lowest weight category to 4.2× in the highest weight category, while female records vary from 4.4× to 3.8×, a weight-adjusted difference of only 10–20%, and an absolute difference of about 30% (i.e., 492 kg vs 348 kg for unlimited weight classes; see Olympic weightlifting records). A study, carried out by analyzing annual world rankings from 1980 to 1996, found that males' running times were, on average, 10% faster than females', with wider hips being a disadvantage for running. However, females have higher endurance than males and there is less than a second difference in the 100m race of Usain Bolt and Florence Griffith Joyner.
In early adolescence, females are on average taller than males (as females tend to go through puberty earlier), but males, on average, surpass them in height in later adolescence and adulthood. In the United States, adult males are on average 9% taller and 16.5% heavier than adult females.
Males typically have larger tracheae and branching bronchi, with about 30 percent greater lung volume per body mass. On average, males have larger hearts, slower heart rates, 10 percent higher red blood cell count, higher hemoglobin, hence greater oxygen-carrying capacity. They also have higher circulating clotting factors (vitamin K, prothrombin and platelets). These differences lead to faster healing of wounds and lower sensitivity to nerve pain after injury. In males, pain-causing injury to the peripheral nerve occurs through the microglia, while in females it occurs through the T cells (except in pregnant women, who follow a male pattern).
Females typically have more white blood cells (stored and circulating), as well as more granulocytes and B and T lymphocytes. Additionally, they produce more antibodies at a faster rate than males, hence they develop fewer infectious diseases and succumb for shorter periods. Ethologists argue that females, interacting with other females and multiple offspring in social groups, have experienced such traits as a selective advantage. Females have a higher sensitivity to pain due to aforementioned nerve differences that increase the sensation, and females thus require higher levels of pain medication after injury. Hormonal changes in females affect pain sensitivity, and pregnant women have the same sensitivity as males. Acute pain tolerance is also more consistent over a lifetime in females than males, despite these hormonal changes. Despite differences in physical feeling, both sexes have similar psychological tolerance to (or ability to cope with and ignore) pain.
In the human brain, a difference between sexes was observed in the transcription of the PCDH11X/Y gene pair unique to Homo sapiens. Sexual differentiation in the human brain from the undifferentiated state is triggered by testosterone from the fetal testis. Testosterone is converted to estrogen in the brain through the action of the enzyme aromatase. Testosterone acts on many brain areas, including the SDN-POA, to create the masculinized brain pattern. The brains of pregnant females carrying male fetuses may be shielded from the masculinizing effects of androgen through the action of sex hormone-binding globulin.
The relationship between sex differences in the brain and human behavior is a subject of controversy in psychology and society at large. Many females tend to have a higher ratio of gray matter in the left hemisphere of the brain in comparison to males. Males on average have larger brains than females; however, when adjusted for total brain volume, the gray matter differences between sexes are almost nonexistent. Thus, the percentage of gray matter appears to be more related to brain size than it is to sex. Differences in brain physiology between sexes do not necessarily relate to differences in intellect. Haier et al. found in a 2004 study that "men and women apparently achieve similar IQ results with different brain regions, suggesting that there is no singular underlying neuroanatomical structure to general intelligence and that different types of brain designs may manifest equivalent intellectual performance". (See the sex and intelligence article for more on this subject.) Strict graph-theoretical analysis of the human brain connections revealed that in numerous graph-theoretical parameters (e.g., minimum bipartition width, edge number, the expander graph property, minimum vertex cover), the structural connectome of women are significantly "better" connected than the connectome of men. It was shown that the graph-theoretical differences are due to the sex and not to the differences in the cerebral volume, by analyzing the data of 36 females and 36 males, where the brain volume of each man in the group was smaller than the brain volume of each woman in the group.
Sexual dimorphism was also described in the gene level and shown to extend from the sex chromosomes. Overall, about 6500 genes have been found to have sex-differential expression in at least one tissue. Many of these genes are not directly associated with reproduction, but rather linked to more general biological features. In addition, it has been shown that genes with sex-specific expression undergo reduced selection efficiency, which leads to higher population frequencies of deleterious mutations and contributes to the prevalence of several human diseases.
Immune function
Sexual dimorphism in immune function is a common pattern in vertebrates and also in a number of invertebrates. Most often, females are more 'immunocompetent' than males. This trait is not consistent among all animals, but differs depending on taxonomy, with the most female-biased immune systems being found in insects. In mammals this results in more frequent and severe infections in males and higher rates of autoimmune disorders in females. One potential cause may be differences in gene expression of immune cells between the sexes. Another explanation is that endocrinological differences between the sexes impact the immune system – for example, testosterone acts as an immunosuppressive agent.
Cells
Phenotypic differences between sexes are evident even in cultured cells from tissues. For example, female muscle-derived stem cells have a better muscle regeneration efficiency than male ones. There are reports of several metabolic differences between male and female cells and they also respond to stress differently. These differences align with concepts of cell autonomous sex identity, where cells maintain intrinsic sex-based traits regardless of systemic influences.
Reproductively advantageous
In theory, larger females are favored by competition for mates, especially in polygamous species. Larger females offer an advantage in fertility, since the physiological demands of reproduction are limiting in females. Hence there is a theoretical expectation that females tend to be larger in species that are monogamous.
Females are larger in many species of insects, many spiders, many fish, many reptiles, owls, birds of prey and certain mammals such as the spotted hyena, and baleen whales such as blue whale. As an example, in some species, females are sedentary, and so males must search for them. Fritz Vollrath and Geoff Parker argue that this difference in behaviour leads to radically different selection pressures on the two sexes, evidently favouring smaller males. Cases where the male is larger than the female have been studied as well, and require alternative explanations.
One example of this type of sexual size dimorphism is the bat Myotis nigricans, (black myotis bat) where females are substantially larger than males in terms of body weight, skull measurement, and forearm length. The interaction between the sexes and the energy needed to produce viable offspring makes it favorable for females to be larger in this species. Females bear the energetic cost of producing eggs, which is much greater than the cost of making sperm by the males. The fecundity advantage hypothesis states that a larger female is able to produce more offspring and give them more favorable conditions to ensure their survival; this is true for most ectotherms. A larger female can provide parental care for a longer time while the offspring matures. The gestation and lactation periods are fairly long in M. nigricans, the females suckling their offspring until they reach nearly adult size. They would not be able to fly and catch prey if they did not compensate for the additional mass of the offspring during this time. Smaller male size may be an adaptation to increase maneuverability and agility, allowing males to compete better with females for food and other resources.
Some species of anglerfish also display extreme sexual dimorphism. Females are more typical in appearance to other fish, whereas males are tiny rudimentary creatures with stunted digestive systems. A male must find a female and fuse with her: he then lives parasitically, becoming little more than a sperm-producing body in what amounts to an effectively hermaphrodite composite organism. A similar situation is found in the Zeus water bug Phoreticovelia disparata where the female has a glandular area on her back that can serve to feed a male, which clings to her (although males can survive away from females, they generally are not free-living). This is taken to the logical extreme in the Rhizocephala crustaceans, like the Sacculina, where the male injects itself into the female's body and becomes nothing more than sperm producing cells, to the point that the superorder used to be mistaken for hermaphroditic.
Some plant species also exhibit dimorphism in which the females are significantly larger than the males, such as in the moss Dicranum and the liverwort Sphaerocarpos. There is some evidence that, in these genera, the dimorphism may be tied to a sex chromosome, or to chemical signalling from females.
Another complicated example of sexual dimorphism is in Vespula squamosa, the southern yellowjacket. In this wasp species, the female workers are the smallest, the male workers are slightly larger, and the female queens are significantly larger than her female workers and male counterparts.
Evolution
In 1871, Charles Darwin advanced the theory of sexual selection, which related sexual dimorphism to sexual selection.
The first step towards sexual dimorphism is the size differentiation of sperm and eggs (anisogamy). Anisogamy and the usually large number of small male gametes relative to the larger female gametes usually lies in the development of strong sperm competition, because small sperm enable organisms to produce a large number of sperm, and make males (or male function of hermaphrodites) more redundant.
Volvocine algae have been useful in understanding the evolution of sexual dimorphism and species like the beetle C. maculatus, where the females are larger than the males, are used to study its underlying genetic mechanisms.
In many non-monogamous species, the benefit to a male's reproductive fitness of mating with multiple females is large, whereas the benefit to a female's reproductive fitness of mating with multiple males is small or nonexistent. In these species, there is a selection pressure for whatever traits enable a male to have more matings. The male may therefore come to have different traits from the female.
These traits could be ones that allow him to fight off other males for control of territory or a harem, such as large size or weapons; or they could be traits that females, for whatever reason, prefer in mates. Male–male competition poses no deep theoretical questions but mate choice does.
Females may choose males that appear strong and healthy, thus likely to possess "good alleles" and give rise to healthy offspring. In some species, however, females seem to choose males with traits that do not improve offspring survival rates, and even traits that reduce it (potentially leading to traits like the peacock's tail). Two hypotheses for explaining this fact are the sexy son hypothesis and the handicap principle.
The sexy son hypothesis states that females may initially choose a trait because it improves the survival of their young, but once this preference has become widespread, females must continue to choose the trait, even if it becomes harmful. Those that do not will have sons that are unattractive to most females (since the preference is widespread) and so receive few matings.
The handicap principle states that a male who survives despite possessing some sort of handicap thus proves that the rest of his genes are "good alleles". If males with "bad alleles" could not survive the handicap, females may evolve to choose males with this sort of handicap; the trait is acting as a hard-to-fake signal of fitness.
| Biology and health sciences | Basics_2 | null |
197191 | https://en.wikipedia.org/wiki/Rock%20dove | Rock dove | The rock dove, rock pigeon, or common pigeon ( also ; Columba livia) is a member of the bird family Columbidae (doves and pigeons). In common usage, it is often simply referred to as the "pigeon", although this is the wild form of the bird; the pigeons most familiar to people are the domesticated form of the wild rock dove.
The domestic pigeon (Columba livia domestica, which includes about 1,000 different breeds) descended from this species. Escaped domestic pigeons have increased the populations of feral pigeons around the world.
Wild rock doves are pale grey with two black bars on each wing, whereas domestic and feral pigeons vary in the colour and pattern of their plumage. Few differences are seen between males and females; i.e. they are not strongly sexually dimorphic. The species is generally monogamous, with two squabs (young) per brood. Both parents care for the young for a time.
Habitats include various open and semi-open environments where they are able to forage on the ground. Cliffs and rock ledges are used for roosting and breeding in the wild. Originating in Southern Europe, North Africa, and Western Asia, pigeons have become established in cities around the world. A commonly cited example of a least-concern species per IUCN among birds, the rock doves are abundant, with an estimated population of 17 to 28 million feral and wild birds in Europe alone and up to 120 million worldwide.
Taxonomy and systematics
The official common name is rock dove, as given by the International Ornithological Congress.
The rock dove was formally described in 1789 by the German naturalist Johann Friedrich Gmelin in his revised and expanded edition of Systema Naturae written by Carl Linnaeus. He placed it with all the other doves and pigeons in the genus Columba and coined the binomial name Columba livia. The genus name Columba is the Latin word meaning "pigeon, dove", whose older etymology comes from the Ancient Greek κόλυμβος (kólumbos), "a diver", hence κολυμβάω (kolumbáō), "dive, plunge headlong, swim". Aristophanes (Birds 304) and others use the word κολυμβίς (kolumbís), "diver", for the name of the bird, because of its swimming motion in the air. The specific epithet livia is a Medieval Latin variant of livida, "livid, bluish-grey"; this was Theodorus Gaza's translation of Greek péleia, "dove", itself thought to be derived from pellós, "dark-coloured". Its closest relative in the genus Columba is the hill pigeon, followed by the other rock pigeons: the snow, speckled, and white-collared pigeons. Pigeon chicks are called "squabs".
Note that members of the lesser known pigeon genus Petrophassa and the speckled pigeon (Columba guinea), also have the common name "rock pigeon".
The rock dove was first described by German naturalist Johann Gmelin in 1789. The rock dove was central to Charles Darwin's discovery of evolution, and featured in four of his works from 1859 to 1872. Darwin posited that, despite wide-ranging morphological differences, the many hundreds of breeds of domestic pigeon could all be traced back to the wild rock dove; in essence human selection of pigeon breeds was analogous to natural selection.
Subspecies
Nine subspecies are recognised:
C. l. livia Gmelin, JF, 1789 – west, central Europe, north Africa to central Asia
C. l. gymnocycla Gray, GR, 1856 – Mauritania and Senegal to south Mali and Ghana
C. l. targia Geyr von Schweppenburg, 1916 – north Mali and south Algeria to central Sudan
C. l. dakhlae Meinertzhagen, R, 1928 – west Egypt
C. l. schimperi Bonaparte, 1854 – east Egypt, south Sudan and Eritrea
C. l. palaestinae Zedlitz, 1912 – Sinai Peninsula (Egypt) to Syria and west, south Arabian Peninsula
C. l. gaddi Zarudny & Loudon, 1906 – east Turkey to Uzbekistan and west, north Afghanistan
C. l. neglecta Hume, 1873 – west Pakistan and east Afghanistan to the Himalayas
C. l. intermedia Strickland, 1844 – south India and Sri Lanka
Description
Centuries of domestication have greatly altered the rock dove. Feral pigeons, which have escaped domestication throughout history, have significant variations in plumage. When not specified, descriptions are for assumed wild type, though the wild type may be on the verge of extinction or already extinct.
The adult of the nominate subspecies of the rock dove is long with a wingspan. Weight for wild or feral rock doves ranges from , though overfed domestic and semidomestic individuals can exceed normal weights. It has a dark bluish-grey head, neck, and chest with glossy yellowish, greenish, and reddish-purple iridescence along its neck and wing feathers. The iris is orange, red, or golden with a paler inner ring, and the bare skin round the eye is bluish-grey. The bill is grey-black with a conspicuous off-white cere, and the feet are purplish-red. Among standard measurements, the wing chord is typically around , the tail is , the bill is around , and the tarsus is .
The adult female is almost identical in outward appearance to the male, but the iridescence on her neck is less intense and more restricted to the rear and sides, whereas that on the breast is often very obscure.
The white lower back of the pure rock dove is its best identification characteristic; the two black bars on its pale grey wings are also distinctive. The tail has a black band on the end, and the outer web of the tail feathers are margined with white. It is strong and quick on the wing, dashing out from sea caves, flying low over the water, its lighter grey rump showing well from above.
Young birds show little lustre and are duller. Eye colour of the pigeon is generally orange, but a few pigeons may have white-grey eyes. The eyelids are orange and encapsulated in a grey-white eye ring. The feet are red to pink.
The subspecies gymnocycla is smaller and very much darker than the nominate subspecies. It is almost blackish on the head, rump and underparts with a white back and the iridescence of the nape extending onto the head. Subspecies targia is slightly smaller than the nominate, with similar plumage, but the back is concolorous with the mantle instead of white. Subspecies dakhlae is smaller and much paler than the nominate. Subspecies schimperi closely resembles targia, but has a distinctly paler mantle. Subspecies palaestinae is slightly larger than schimperi and has darker plumage. Subspecies gaddi is larger and paler than palaestinae, with which it intergrades in the west. Subspecies neglecta it is similar to the nominate in size but darker, with a stronger and more extensive iridescent sheen on the neck. It intergrades with gaddi in the south. Subspecies intermedia is similar to neglecta but darker, with a less contrasting back.
There have been numerous skeletal descriptions of the rock dove and the associated muscles including those of the eye, jaw, neck, and throat. The skull is dominated by the rostrum, eye socket, and braincase. The quadrate bone is relatively small and mobile and connects the rest of the cranium to the lower jaw. The latter has an angled shape in side view because the long-axis of the front half of the lower jaw is at a 30° angle to the back half. Beneath the skull, the hyoid skeleton involves three mid-line structures and a pair of elongate structures that stem from between the junction of the back two structures. The anterior structure (the paraglossum or entoglossum) is unpaired and shaped like an arrowhead.
When circling overhead, the white underwing of the bird becomes conspicuous. In its flight, behaviour, and voice, which is more of a dovecot coo than the phrase of the wood pigeon, it is a typical pigeon. Although it is a relatively strong flier, it also glides frequently, holding its wings in a very pronounced V shape as it does. As prey birds, they must keep their vigilance, and when disturbed a pigeon within a flock will take off with a noisy clapping sound that cues for other pigeons to take to flight. The noise of the take-off increases the faster a pigeon beats its wings, thus advertising the magnitude of a perceived threat to its flockmates.
Feral pigeons are essentially the same size and shape as the original wild rock dove, but often display far greater variation in colour and pattern compared to their wild ancestors. The blue-barred pattern which the original wild rock dove displays is generally less common in more urban areas. Urban pigeons tend to have darker plumage than those in more rural areas.
Pigeons feathers have two types of melanin (pigment) – eumelanin and pheomelanin. A study of melanin in the feathers of both wild rock and domestic pigeons, of different coloration types and known genetic background, measured the concentration, distribution and proportions of eumelanin and pheomelanin and found that gene mutations affecting the distribution, amounts and proportions of pigments accounted for the greater variation of coloration in domesticated birds than in their wild relations. Eumelanin generally causes grey or black coloration, while pheomelanin results in a reddish-brown colour. Other shades of brown may be produced through different combinations and concentrations of the two colours. As in other animals, white pigeons have little to no pigment. Darker birds may be better able to store trace metals in their feathers due to their higher concentrations of melanin, which may help mitigate the negative effects of the metals, the concentrations of which are typically higher in urban areas. Pigeons, especially homing or carrier breeds, are well known for their ability to find their way home from long distances. Despite these demonstrated abilities, wild rock doves are sedentary and rarely leave their local areas. It is hypothesized that in their natural, arid habitat, they rely on this sense to navigate back home after foraging as deserts rarely possess navigational landmarks that may be used.
A rock dove's lifespan ranges from 3–5 years in the wild to 15 years in captivity, though longer-lived specimens have been reported. The main causes of mortality in the wild are predators and persecution by humans. Some sources state the species was first introduced to North America in 1606 at Port Royal, Nova Scotia. Although other sources cite Plymouth and Jamestown settlements in the early 17th century as the first place for species introduction in North America.
Vocalizations
The call is a soft, slightly wavering, coo. Ornithologist David Sibley describes the display call as a whoo, hoo-witoo-hoo, whereas the Cornell Lab of Ornithology describes it as a Coo, roo-c'too-coo. Variations include an alarm call, a nest call, and noises made by juveniles. Sibley describes the nest call as a repeated hu-hu-hurrr. When displaying, songs are partly sexual, partly threatening. They are accompanied by an inflated throat, tail fanning, strutting, and bowing. The alarm call, given at sight of predators, is a grunt-like oorhh.
Non-vocal sounds include a loud flapping noise at take-off, feet stomping, hisses, and beak snapping. Wings may also be clapped during flights, usually during display fights or after copulation. Juveniles particularly snap their bills, usually to respond to nest invasion. The foot stomping appears deliberate, though for what purpose is unclear. Foot stomping is done with a certain foot first, showing that rock doves have "footedness", similar to human handedness.
Osmoregulation
Distribution and habitat
Before the Columbian Exchange, rock doves were restricted to a natural resident range in western and southern Europe, North Africa, and extending into South Asia. They were carried into the New World aboard European ships between 1603 and 1607. The species (including ferals) has a large range, with an estimated global extent of occurrence of . It has a large global population, including an estimated 17 to 28 million individuals in Europe. Fossil evidence suggests the rock dove originated in southern Asia, and skeletal remains, unearthed in Israel, confirm its existence there for at least 300,000 years. However, this species has such a long history with humans that it is impossible to identify its original range exactly.
Wild pigeons reside in rock formations and cliff faces, settling in crevices to nest. They nest communally, often forming large colonies of many hundreds of individuals. Wild nesting sites include caves, canyons, and sea cliffs. They will even live in the Sahara so long as an area has rocks, water, and some plant matter. They prefer to avoid dense vegetation.
Rock doves have a commensal relationship with humans, gaining both ample access to food and nesting spots in civilized areas. Human structures provide an excellent imitation of cliff structures, making rock doves very common around human habitation. Skyscrapers, highway overpasses, farm buildings, abandoned buildings, and other human structures with ample crevices are conducive to rock dove nesting. Thus the modern range of the rock dove is due in large part to humans. Agricultural settlements are favoured over forested ones. Ideal human nesting attributes combine areas with tall buildings, green spaces, ample access to human food, and schools. Conversely, suburban areas which are far from city centers and have high street density are the least conducive to pigeons. Their versatility among human structures is evidenced by a population living inside a deep well in Tunisia.
Feral pigeons are usually unable to find these accommodations, so they must nest on building ledges, walls or statues. They may damage these structures via their faeces; starving birds can only excrete urates, which over time corrodes masonry and metal. In contrast, a well-fed bird passes mostly solid faeces, containing only small amounts of uric acid.
Behaviour and ecology
Pigeons are often found in pairs in the breeding season, but are usually gregarious.
Breeding
The rock dove breeds at any time of the year, but peak times are spring and summer. Nesting sites are along coastal cliff faces, as well as the artificial cliff faces created by apartment buildings with accessible ledges or roof spaces. Pigeons can compete with native birds for nest sites. For some avian species, such as seabirds, it could be a conservation issue. Current evidence suggests that wild, domestic and feral pigeons mate for life, although their long-term bonds are not unbreakable. They are socially monogamous, but extra-pair matings do occur, often initiated by males. Due to their ability to produce crop milk, pigeons can breed at any time of year. Pigeons breed when the food supply is abundant enough to support embryonic egg development, which in cities, can be any time of the year. Laying of eggs can take place up to six times per year.
Pigeons are often found in pairs during the breeding season, but usually the pigeons are gregarious, living in flocks of 50 to 500 birds (dependent on the food supply).
Courtship rituals can be observed in urban parks at any time of the year. The male on the ground or rooftops puffs up the feathers on his neck to appear larger and thereby impress or attract attention. He approaches the female at a rapid walking pace while emitting repetitive quiet notes, often bowing and turning as he comes closer. At first, the female invariably walks or flies a short distance away and the male follows her until she stops. At this point, he continues the bowing motion and very often makes full- or half-pirouettes in front of the female. The male then proceeds to feed the female by regurgitating food, as they do when feeding the young. The male then mounts the female, rearing backwards to be able to join their cloacae. The mating is very brief, with the male flapping his wings to maintain balance on top of the female.
The nest is a flimsy platform of straw and sticks, laid on a ledge, under cover, often on the window ledges of buildings. Two white eggs are laid; incubation, shared by both parents, lasts 17 to 19 days. The newly hatched squab (nestling) has pale yellow down and a flesh-coloured bill with a dark band. For the first few days, the baby squabs are tended and fed (through regurgitation) exclusively on "crop milk" (also called "pigeon milk" or "pigeon's milk"). The pigeon milk is produced in the crops of both parents in all species of pigeon and dove. Pigeons are altricial and their fledging period is about 30 days.
Feeding
Rock doves are omnivorous, but prefer plant matter: chiefly fruits and grains.
Studies of pigeons in a semi-rural part of Kansas found that their diet includes the following: 92% maize, 3.2% oats, 3.7% cherry, along with small amounts of knotweed, elm, poison ivy and barley. Feral pigeons can be seen eating grass seeds and berries in parks and gardens in the spring, but plentiful sources exist throughout the year from scavenging (e.g., food remnants left inside of dropped fast food cartons, in the form of popcorn, cake, peanuts, bread and currants) and they also eat insects and spiders. Additional food is also usually available from waste bins, tourists or residents who feed bird seed to pigeons for reasons such as empathy, fun, tradition and as a means for social interaction. Pigeons tend to congregate in large, often thick flocks when feeding on discarded food, and may be observed flying skillfully around trees, buildings, telephone poles and cables and even through moving traffic just to reach a food source.
Pigeons feed on the ground in flocks or individually. Pigeons are naturally granivorous, eating seeds that fit down their gullet. They may sometimes consume small invertebrates such as worms or insect larvae as a protein supplement. As they do not possess an enlarged cecum as in European wood pigeons, they cannot digest adult plant tissue; the various seeds they eat contain the appropriate nutrients they require. While most birds take small sips and tilt their heads backwards when drinking, pigeons are able to dip their bills into the water and drink continuously, without having to tilt their heads back. In cities they typically resort to scavenging human garbage, as unprocessed grain may be impossible to find. Pigeon groups typically consist of producers, which locate and obtain food, and scroungers, which feed on food obtained by the producers. Generally, groups of pigeons contain a greater proportion of scroungers than producers.
Preening
Pigeons primarily use powder down feathers for preening, which gives a soft and silky feel to their plumage. They have no preen gland or at times have very rudimentary preen glands, so oil is not used for preening. Rather, powder down feathers are spread across the body. These have a tendency to disintegrate, and the powder, akin to talcum powder, helps maintain the plumage. Some varieties of domestic pigeon have modified feathers called "fat quills". These feathers contain yellow, oil-like fat that derives from the same cells as powder down. This is used while preening and helps reduce bacterial degradation of feathers by feather bacilli.
Survival
Predators
With only their flying abilities protecting them from predation, rock pigeons are a favourite almost around the world for a wide range of raptors. In fact, with feral pigeons existing in almost every city in the world, they may form the majority of prey for several raptor species that live in urban areas. Peregrine falcons and Eurasian sparrowhawks are natural predators of pigeons and quite adept at catching and feeding upon this species. Up to 80% of the diet of peregrine falcons in several cities that have breeding falcons is composed of feral pigeons. Some common predators of feral pigeons in North America are raccoons, opossums, red-tailed hawks, great horned owls, eastern screech owls, and accipiters. The birds that prey on pigeons in North America can range in size from American kestrels to golden eagles and may even include crows, gulls and ravens. On the ground the adults, their young and their eggs are at risk from feral and domestic cats. Doves and pigeons are considered to be game birds, since many species are hunted and used for food in many of the countries in which they are native.
The body feathers have dense, fluffy bases and are loosely attached to the skin, hence they drop out easily. When a predator catches a pigeon large numbers of feathers come out in the attacker's mouth and the pigeon may use this temporary distraction to make an escape. It also tends to drop the tail feathers when preyed upon or under traumatic conditions, probably as a distraction mechanism.
Parasites
Pigeons may harbour a diverse parasite fauna.
They often host the intestinal helminths Capillaria columbae and Ascaridia columbae. Their ectoparasites include the ischnoceran lice Columbicola columbae, Campanulotes bidentatus compar, the amblyceran lice Bonomiella columbae, Hohorstiella lata, Colpocephalum turbinatum, the mites Tinaminyssus melloi, Dermanyssus gallinae, Dermoglyphus columbae, Falculifer rostratus and Diplaegidia columbae. The hippoboscid fly Pseudolynchia canariensis is a typical blood-sucking ectoparasite of pigeons in tropical and subtropical regions.
Relationship to humans
Domestication
Rock doves have been domesticated for several thousand years, giving rise to the domestic pigeon (Columba livia domestica). They may have been domesticated as long as 5,000 years ago. Numerous breeds of fancy pigeons of all sizes, colours, and types have been bred. Domesticated pigeons are used as homing pigeons as well as food and pets. They were in the past also used as carrier pigeons.
War pigeons
So-called war pigeons have played significant roles during wartime, and many pigeons have received awards and medals for their services in saving hundreds of human lives.
Medical uses
Pigeons have notably been "employed" as medical imaging data sorters. They have been successfully trained under research conditions to examine data on a screen for the purposes of detecting breast cancer. They appear to use their innate visual navigation skills to do so.
Feral pigeon
Many domestic birds have gotten lost, escaped or been released over the years and have given rise to feral pigeons. These show a variety of plumages, although many have the blue-barred pattern as does the pure rock dove. Feral pigeons are found in cities and towns all over the world. The scarcity of the pure wild species is partly due to interbreeding with feral birds.
Human health
Contact with pigeon droppings poses a minor risk of contracting histoplasmosis, cryptococcosis and psittacosis, and long-term exposure to both droppings and feathers can induce an allergy known as bird fancier's lung.
Pigeons are not a major concern in the spread of West Nile virus: though they can contract it, they apparently do not transmit it.
Some contagions are transmitted by pigeons; for example, the bacteria Chlamydophila psittaci is endemic among pigeons and causes psittacosis in humans. It is generally transmitted from handling pigeons or their droppings (more commonly the latter). Psittacosis is a serious disease but rarely fatal (less than 1%). Pigeons are also important vectors for various species of the bacteria Salmonella, which causes diseases such as salmonellosis and paratyphoid fever.
Pigeons are also known to host avian mites, which can infest human habitation and bite humans, a condition known as gamasoidosis. However, infesting mammals is relatively rare.
Avian influenza
Pigeons may, however, carry and spread avian influenza. One study has shown that adult pigeons are not clinically susceptible to the most dangerous strain of avian influenza, H5N1, and that they do not transmit the virus to poultry. Other studies have presented evidence of clinical signs and neurological lesions resulting from infection but found that the pigeons did not transmit the disease to poultry reared in direct contact with them. Pigeons were found to be "resistant or minimally susceptible" to other strains of avian influenza, such as the H7N7.
Research into whether pigeons play a part in spreading bird flu have shown pigeons do not carry the deadly H5N1 strain. Three studies have been done since the late 1990s by the US Agriculture Department's Southeast Poultry Research Laboratory in Athens, Georgia, according to the center's director, David Swayne. The lab has been working on bird flu since the 1970s. In one experiment, researchers squirted into pigeons' mouths liquid drops that contained the highly pathogenic H5N1 virus from a Hong Kong sample. The birds received 100 to 1,000 times the concentration that wild birds would encounter in nature. "We couldn't infect the pigeons", Swayne said. "So that's good news."
Stages of lifecycle
| Biology and health sciences | Columbiformes | null |
197301 | https://en.wikipedia.org/wiki/Chassepot | Chassepot | The Chassepot (pronounced SHAS-poh), officially known as Fusil modèle 1866, was a bolt-action military breechloading rifle. It is famous for having been the arm of the French forces in the Franco-Prussian War of 1870–1871. It replaced an assortment of Minié muzzleloading rifles, many of which were converted in 1864 to breech loading (the Tabatière rifles). An improvement to existing military rifles in 1866, the Chassepot marked the commencement of the era of modern bolt action, breech-loading military rifles. The Gras rifle was an adaption of the Chassepot designed to fire metallic cartridges introduced in 1874.
It was manufactured by Manufacture d'armes de Saint-Étienne (MAS), Manufacture d'Armes de Châtellerault (MAC), Manufacture d'Armes de Tulle (MAT) and, until 1870, in the Manufacture d'Armes de Mutzig in the former Château des Rohan. Many were also manufactured under contract in England (the "Potts et Hunts" Chassepots delivered to the French Navy), in Belgium (Liege), and in Italy at Brescia (by Glisenti). The approximate number of Chassepot rifles available to the French Army in July 1870 was 1,037,555 units. Additionally, state manufacturies could deliver 30,000 new rifles monthly. Gun manufacturers in England and Austria also produced Chassepot rifles to support the French war effort. The Josef und Franz Werndl & Co. in Steyr, Austria delivered 12,000 Chassepot carbines and 100,000 parts to France in 1871. Manufacturing of the Chassepot rifle ended in February 1875, four years after the end of the Franco-Prussian War, with approximately 700,000 more Chassepot rifles made between September 1871 and July 1874.
History
The Chassepot was named after its inventor, Antoine Alphonse Chassepot (1833–1905), who, from the mid-1850s onwards, had constructed various experimental forms of breech loaders. The first two models of the Chassepot still used percussion cap ignition. The third model, using a similar system to the Prussian Dreyse needle gun, became the French service weapon on 30 August 1866. In the following year it made its first appearance at the Battle of Mentana on 3 November 1867, where it inflicted severe losses upon Giuseppe Garibaldi's troops. It was reported at the French Parliament that "Les Chassepots ont fait merveille!", ("The Chassepots have done wonderfully!") The heavy cylindrical lead bullets fired at high velocity by the Chassepot rifle inflicted wounds that were even worse than those of the Minié rifle. By 1868, the entire French active army had been re-armed with the Chassepot.
In the Franco-Prussian War (1870–1871), the Chassepot met its Prussian counterpart, the Dreyse needle-fire rifle. The Chassepot had several advantages over the Dreyse. It featured a rubber obturator on its bolt head to provide a more efficient gas-seal. Although it fired a smaller caliber (11 mm vs. 15.4 for the Dreyse), the Chassepot ammunition had more gunpowder (5.68 grams vs 4.85 grams), resulting in higher muzzle velocity (436 meters per second, 33% over the Dreyse), a flatter trajectory and a longer range. Thus the sights on the Chassepot could be elevated up to 1,600 meters, while the maximum sight setting of the Dreyse was only 600 meters. The Chassepot had a weight of 4.1 kg versus 4.57 kg for the needle-fire rifle. It was also shorter (1310 mm vs. 1424 mm).
After the war, 20,000 captured Chassepot rifles were sold to the Shah of the Persian Qajar dynasty.
In 1872 the Empire of Brazil purchased 8631 Chassepots; after being faced with a possible war involving Argentine claims over Paraguay. The weapons, however, were never officially distributed to the Army, since the decision to buy the Comblain had already been made and because of issue with the cartridge's reliability. The weapons ended their career in deposit or were handed over to police forces and shooting clubs. Some of these weapons were possibly used by rebels during the War of Canudos, they may have been captured from the Bahia police after the engagement at Maceté. Chassepots were used during the Federalist Revolution and by the rebels of the 1923 revolution in Rio Grande do Sul.
In June 1880, some 40 1866 Chassepots with a few bayonets were delivered to the port of Buenos Aires. Those Chassepots were of 11mm calibre and were possibly rechambered for the Gras cartridge (They were delivered together with Gras rifles in a shipment of 450 weapons).
Surplus Chassepot were exported to China. Some of the warriors of the Ethiopian Empire were equipped with Chassepot rifles during the first Italo-Ethiopian War of 1896.
Technology
Bolt mechanism
The breech was closed by a bolt similar to those of more modern rifles. Amongst the technical features of interest introduced in 1866 on the Chassepot rifle was the method of obturation of the bolt by a segmented rubber ring which expanded under gas pressure and thus sealed the breech when the shot was fired. This simple yet effective technology was successfully adapted to artillery in 1877 by Colonel de Bange, who invented grease-impregnated asbestos pads to seal the breech of his new cannons (the De Bange system).
Cartridge
The Chassepot used a paper cartridge that many refer to as 'combustible', although in reality it was quite the opposite. It held an round-headed cylindro-conoidal lead bullet that was wax paper patched. An inverted standard percussion cap was at the rear of the paper cartridge and hidden inside. It was fired by the Chassepot's needle (a sharply pointed firing pin) upon pressing the trigger. While the Chassepot's ballistic performance and firing rates were excellent for the time, burnt paper residues as well as black powder fouling accumulated in the chamber and bolt mechanism after continuous firing. The bolt's rubber obturator eroded in action but was easily replaced in the field by infantrymen. The older Dreyse needle gun and its cartridge had been designed to minimize those problems but to the detriment of its ballistic properties.
To correct this problem the Chassepot was replaced in 1874 by the Gras rifle which used a centerfire drawn brass metallic cartridge. Otherwise, the Gras rifle was basically identical in outward appearance to the Chassepot rifle. Nearly all rifles of the older Chassepot model (Mle 1866) remaining in store were eventually converted to take the 11 mm Gras metallic cartridge ammunition (fusil Modèle 1866/74). About 665,327 Chassepot rifles had been captured by the German coalition that defeated France in 1871. Large numbers of these captured Chassepot rifles were shortened and converted to 11 mm Mauser metallic cartridge. It served with cavalry units of the Kingdom of Saxony and of the Kingdom of Bavaria. Others were disposed of "as is" with British surplus dealers. In most but not all, the French receiver markings on these German-captured Chassepot rifles had been erased.
Gallery
| Technology | Specific firearms | null |
197341 | https://en.wikipedia.org/wiki/Arleigh%20Burke-class%20destroyer | Arleigh Burke-class destroyer | The Arleigh Burke class of guided-missile destroyers (DDGs) is a United States Navy class of destroyer centered around the Aegis Combat System and the SPY-1D multi-function passive electronically scanned array radar. The class is named after Admiral Arleigh Burke, an American destroyer officer in World War II and later Chief of Naval Operations. With an overall length of , displacement ranging from 8,300 to 9,700 tons, and weaponry including over 90 missiles, the Arleigh Burke-class destroyers are larger and more heavily armed than many previous classes of guided-missile cruisers.
These warships are multi-mission destroyers able to conduct anti-aircraft warfare with Aegis and surface-to-air missiles; tactical land strikes with Tomahawk missiles; anti-submarine warfare (ASW) with towed array sonar, anti-submarine rockets, and ASW helicopters; and anti-surface warfare (ASuW) with ship-to-ship missiles and guns. With upgrades to their AN/SPY-1 radar systems and their associated missile payloads as part of the Aegis Ballistic Missile Defense System, as well as the introduction of the AN/SPY-6 radar system, the class has also evolved capability as mobile anti-ballistic missile and anti-satellite platforms.
The lead ship of the class, , was commissioned during Admiral Burke's lifetime on 4 July 1991. With the decommissioning of the last , , on 2005, the Arleigh Burke-class ships became the U.S. Navy's only active destroyers until the became active in 2016. The Arleigh Burke class has the longest production run of any U.S. Navy surface combatant. all seventy-four built are active, with nineteen more planned to enter service.
Characteristics
Variants
The Arleigh Burke-class destroyer has four variants, referred to as "Flights". Newer Flights incorporate technological advancements.
Flight I: DDGs 51–71
Flight II: DDGs 72–78
Flight IIA: DDGs 79–124 and DDG 127
Flight III: DDGs 125–126 and DDG 128 onwards
Structure
The Arleigh Burke-class ships are among the largest destroyers built in the United States; only the , (), and classes () are longer. The Arleigh Burke class was designed with a new large, water-plane area-hull form characterized by a wide flaring bow, which significantly improves seakeeping ability and permits high speed in high sea states. The class's design incorporates stealth techniques, such as the angled (rather than traditional vertical) surfaces and the raked tripod mainmast, which make the ship more difficult to detect by radar.
Its designers incorporated lessons from the , which the Navy deemed too expensive to continue building and difficult to upgrade further. For these destroyers, the U.S. Navy returned to all-steel construction, except the mast made of aluminum. The Ticonderogas had combined a steel hull with a superstructure made of lighter aluminum to reduce top weight, but the lighter metal proved vulnerable to cracking. Aluminum is also less fire-resistant than steel; a 1975 fire aboard gutted her aluminum superstructure. Battle damage to Royal Navy ships exacerbated by their aluminum superstructures during the 1982 Falklands War supported the decision to use steel. Other lessons from the Falklands War led to the Navy's decision to protect the Arleigh Burke class's vital spaces with double-spaced steel layers, which create a buffer against anti-ship missiles (AShMs), and Kevlar spall liners.
Passive defenses
Arleigh Burke destroyers are equipped with AN/SLQ-32 electronic warfare (EW) suites that provide electronic support. Vessels with the SLQ-32(V)3, SLQ-32(V)6, or SLQ-32(V)7 variant can jam radars.
The destroyers have Mark 36 infrared and chaff decoy launchers, as well as Nulka decoy launchers, for spoofing incoming AShMs. For defeating incoming torpedoes, the class has two AN/SLQ-25 Nixie towed countermeasures. The ships' Prairie-Maskers can reduce their radiated noise.
A collective protection system makes the Arleigh Burke class the first U.S. warships designed with an air-filtration system against nuclear, biological, and chemical warfare (NBC). Other NBC defenses include double air-locked hatches, pressurized compartments, and an external countermeasure washdown system. The class's electronics are hardened against electromagnetic pulses. Fire suppression equipment includes water sprinklers in the living quarters and combat information center (CIC).
Weapon systems
The Arleigh Burke class are multi-mission ships with numerous combat systems, including anti-aircraft missiles, land attack missiles, ship-to-ship missiles, and an anti-submarine warfare (ASW) system. Missiles are stored in and fired from Mark 41 Vertical Launching System (VLS) cells; with 90 cells on Flights I–II and 96 cells starting with Flight IIA, the Arleigh Burkes are more heavily armed than many preceding guided-missile cruiser classes. The Arleigh Burke-class destroyer is equipped with the Aegis Combat System, which combines information from the ship's sensors to display a coherent image of the environment and guides weapons to targets using advanced tracking and fire control.
Their main radar differs from traditional mechanically rotating radars. Instead, Aegis uses the AN/SPY-1D passive electronically scanned array (or the AN/SPY-6 active electronically scanned array on Flight III ships), which allows continual tracking of targets simultaneous to area scans. The system's computer control also allows centralization of the previously separate tracking and targeting functions. The system is resistant to electronic countermeasures.
The Standard Missile SM-2MR/ER and SM-6 provide area air defense, though they may also be used in a secondary anti-ship role. The SM-2 uses semi-active radar homing (SARH); up to three targets may be simultaneously intercepted as the Arleigh Burkes have three AN/SPG-62 fire-control radars for terminal target illumination. The SM-6, which provides over-the-horizon defense, and the SM-2 Block IIIC feature a dual-mode seeker with active radar homing (ARH) capability; they do not have to rely on external illumination, so more targets may be intercepted simultaneously.
Flights IIA and III—and modernized Flight I and II ships—can carry RIM-162 Evolved SeaSparrow Missiles (ESSMs), which provide medium-range air defense and are also capable of targeting other ships. ESSM is small enough to be quad-packed into a single Mk 41 VLS cell. ESSM Block 1 uses SARH, guided similarly to older SM-2s. ESSM Block 2, which achieved initial operating capability (IOC) in 2021, features a dual-mode seeker with ARH capability.
The SM-3, SM-6, and SM-2ER Block IV provide Ballistic Missile Defense (BMD), the SM-3 being an exoatmospheric interceptor and the latter two having terminal phase anti-ballistic capability. So vital has the Aegis BMD role become that all ships of the class are being updated with BMD capability. By January 2023, there were 51 BMD-capable Arleigh Burke-class destroyers. Flight III ships have been delivered since 2023 with AN/SPY-6(V)1 radars and improved BMD capabilities; Flight IIA ships are also planned to receive these upgrades with AN/SPY-6(V)4 radar retrofits.
Flights I and II carry two stand-alone Harpoon anti-ship missile launchers for a total of four or eight Harpoons, providing an anti-ship capability with a range in excess of . During Exercise RIMPAC 2024, DDG-62, a Flight I ship, launched a Naval Strike Missile (NSM); the launchers for the Harpoons were removed to make room for the NSM's proprietary launch boxes.
The class can perform tactical land strikes with VLS-launched Tomahawks. With the development of the Tomahawk Block V, all existing Block IV Tomahawks carried will be converted to the Block V. The Tomahawk Block Va version is called the Maritime Strike version, and it provides anti-ship capability in addition to its land attack role. The Block Vb version features the Joint Multi-Effects Warhead System for hitting a wider variety of land targets.
Arleigh Burke-class ships have the AN/SQQ-89 ASW combat system, which is integrated with Aegis. It encompasses the AN/SQS-53C bow-mounted sonar and a towed array sonar, though several Flight IIA ships do not have a towed array. The towed array is either the AN/SQR-19 Tactical Towed Array Sonar (TACTAS) or the newer TB-37U Multi-Function Towed Array (MFTA). The ships can carry standoff RUM-139 vertical launch anti-submarine rockets. A Mark 32 triple torpedo tubes mount on each side of the ship can fire Mark 46, Mark 50, or Mark 54 lightweight torpedoes for short-range ASW. The ships can detect anti-ship mines at a range of about 1,400 meters.
All ships of the class are fitted with at least one Phalanx close-in weapon system (CIWS), which provides point defense against air and surface threats. Eight ships (DDG 51, DDG 64, DDG 71, DDG 75, DDG 78, DDG 80, DDG 84, DDG 117) are equipped with one SeaRAM CIWS for improved self-defense. Arleigh Burkes can also carry two 25 mm Mk 38 machine gun systems, one on each side of the ship, designed to counter fast surface craft. There are numerous mounts for crew-served weapons like the M2 Browning.
Located on the forward deck is the 5-inch (127 mm) Mark 45 gun. Directed by the Mark 34 Gun Weapon System, it can be used in anti-ship, anti-air, and naval gunfire support (NGFS) roles. It can fire 16–20 rounds per minute and has a range of . Arleigh Burkes can stow 680 5-inch rounds.
As of 2023, six destroyers (DDG 100, DDG 104, DDG 105, DDG 106, DDG 111, DDG 113) are equipped with the Optical Dazzling Interdictor, Navy (ODIN), a directed energy weapon that can target unmanned vehicles. DDG 88 is equipped with the higher-power High Energy Laser with Integrated Optical-dazzler and Surveillance (HELIOS).
Aircraft
Flights IIA and III have two hangars for stowing MH-60 helicopters. Their Light Airborne Multi-Purpose System (LAMPS) helicopter system improves the ship's capabilities by enabling the MH-60 to monitor submarines and surface ships, launch torpedoes and missiles against them, and provide fire support during insertions/extractions with machine guns and Hellfire anti-armor guided missiles. The helicopters also serve in a utility role, able to perform vertical replenishment, search and rescue, medical evacuation, communications relay, and naval gunfire spotting and controlling.
In March 2022, an Arleigh Burke destroyer was deployed with an AAI Aerosonde unmanned aerial vehicle (UAV). The aircraft is under demonstration for Flight I and II ships, which do not have accommodations for permanently storing helicopters. The Aerosonde has a small enough footprint to be stowed on those destroyers. It can perform missions such as intelligence, surveillance, and reconnaissance at a much lower cost than manned helicopters.
Development
Origins and Flight I
The Chief of Naval Operations (CNO) from 1970 to 1974, Admiral Elmo Zumwalt, sought to improve the U.S. Navy through modernization at minimal cost. Zumwalt's approach to the fleet was a "high-low mix"—a few high-end, high-cost warships supplemented by numerous low-end, low-cost warships. The introduction of the Aegis-equipped Ticonderoga-class cruiser in the early 1980s filled the high end. The Navy started work to develop a lower-cost Aegis-equipped vessel to fill the low end and replace the aging destroyers.
In 1980, the U.S. Navy initiated design studies with seven contractors. By 1983, the number of competitors had been reduced to three: Bath Iron Works, Ingalls Shipbuilding, and Todd Shipyards. On 3 April 1985, Bath Iron Works received a US$321.9 million contract to build the first of the class, USS Arleigh Burke. Gibbs & Cox was awarded the contract to be the lead ship design agent. The Navy contracted Ingalls Shipbuilding to build the second ship.
Political restraints led to design restrictions, including the absence of helicopter hangars, a displacement limit of 8,300 tons, and a 50-foot shorter hull than the Ticonderoga's. The designers were forced to make compromises, such as a wide flaring bow. To compensate for the limited length, the originally-planned 80,000 shaft horsepower (shp) LM2500 gas turbines were upgraded to 100,000 shp. No main gun was included in the original design, later amended to include an OTO Melara 76 mm, before finally selecting the 5-inch/54-caliber Mark 45. Despite their constraints, the designers benefitted from insight gained from previous classes; for example, they chose an all-steel superstructure to improve survivability.
The total cost of the first ship was $1.1 billion, the other $778 million being for the ship's weapons systems. USS Arleigh Burke was laid down by the Bath Iron Works at Bath, Maine, on 6 December 1988, and launched on 16 September 1989 by Mrs. Arleigh Burke. The Admiral himself was present at her commissioning ceremony on 4 July 1991, held on the waterfront in downtown Norfolk, Virginia. Orders for Flight I ships continued through 1995.
Flight II
The Flight II iteration of the class was introduced in FY1992. The incorporation of the AN/SRS-1A(V) Combat Direction Finding enhanced detection of signals. The TADIX-B, JTIDS Command and Control Processor, and Link 16 improved communication with other assets. The SLQ-32 EW suite was upgraded to (V)3, and the SPS-67(V)3 surface search radar was upgraded to (V)5. Flight II also gained the capability to launch and control the SM-2ER Block IV. An expansion of fuel capacity slightly increased the displacement.
Flight IIA
The Flight IIA design was first procured in FY1994. Among the additions are two hangars and support facilities for ASW helicopters, Cooperative Engagement Capability (CEC), the Kingfisher mine detection system, and five blast-resistant bulkheads. To accommodate the hangars, the length was increased to , and the rear-facing SPY-1D arrays are mounted one deck (eight feet) higher to prevent a blind spot. Flight IIA replaced retractable missile loading cranes on the forward and aft VLS with a total of six additional cells. The propellers are of a different design to reduce cavitation. New fiber optics improved bandwidth and helped reduce weight gain. Systems removed from Flight IIA include the Harpoon missile launchers and, starting with , the forward Phalanx CIWS. Flight IIA ships were initially built without the AN/SQR-19 TACTAS, though later units were subsequently installed with TACTAS.
Starting with , the longer 5-inch/62-caliber (127 mm) Mark 45 Mod 4 gun was installed. Later Flight IIA ships starting with use the BridgeMaster E as their navigation radar instead of the AN/SPS-73(V)12. Subsequent Flight IIA ships employ additional signature-reduction measures: the hangars of DDG 86 onwards are made of composite materials, and the exhaust funnels of DDG 89 onwards are shrouded by the superstructure. The use of the improved SPY-1D(V) radar, starting with , enhances the ships' ability to filter out clutter and resist electronic attack.
Several Flight IIA ships were constructed without any Phalanx CIWS because of the planned Evolved SeaSparrow Missile; the Navy had initially decided that ESSM made Phalanx redundant. However, the Navy later changed its mind and decided to retrofit all IIA ships to carry at least one Phalanx CIWS by 2013.
DDGs 91–96 (USS Pinckney, , , , , and ) were built with superstructure differences to accommodate the AN/WLD-1 Remote Minehunting System (RMS). However, only Pinckney, Momsen, and Bainbridge were installed with the system before the RMS program was canceled.
Modernization
Efforts to modernize the Arleigh Burke class began amid congressional concerns over the retirement of the . In 1996, the Navy began a program to field the Extended Range Guided Munition (ERGM) for the DDG 51 class. The ERGM was to extend the class's 5-inch Mark 45 gun range to . It necessitated a modification of the gun; the 62-caliber Mark 45 Mod 4 was created and installed on DDG 81 and onwards in anticipation of the ERGM. However, the ERGM was canceled in 2008.
The current DDG 51 modernization program is designed to provide mid-life upgrades to ensure the destroyers remain effective with service lives of at least 35 years. Modernization of existing ships provides commonality with in-production ships. The program's goals are reduced manning, increased mission effectiveness, and reduced total cost. Mid-life modernization of Flight I and II ships is done in two phases: the first phase updates the hull, mechanical, and electrical (HM&E) systems, while the second phase focuses on Aegis Combat System upgrades and introduces an Open Architecture Computing Environment (OACE). By 2017, modernization technologies were introduced to production ships, and the Navy started modernization of Flight IIA ships through a single process combining both phases of upgrading. The capabilities of modernized destroyers include CEC, Integrated Air and Missile Defense (IAMD), ESSM support, improved electronic support with Surface Electronic Warfare Improvement Program (SEWIP) Block 2, improved data processing with Boeing's Gigabit Ethernet Data Multiplex System, and improvements to littoral warfare.
In July 2010, BAE Systems announced it had been awarded a contract to modernize 11 ships. In May 2014, USNI News reported that 21 of the 28 Flight I and II Arleigh Burke-class destroyers would not receive the full mid-life upgrade that included electronics and Aegis Baseline 9 software for SM-6 compatibility; instead, they would retain the basic BMD 3.6.1 software in a $170 million upgrade concentrating on HM&E systems, and on some ships, their anti-submarine suite. Seven Flight I ships—DDGs 51–53, 57, 61, 65, 69—received the full $270 million Baseline 9 upgrade. Deputy of surface warfare Dave McFarland said that this change was due to the budget cuts in the Budget Control Act of 2011.
In 2016, the Navy announced it would begin outfitting 34 Flight IIA Arleigh Burkes with a hybrid-electric drive (HED) to lower fuel costs. The four LM2500 gas turbines of the class are most efficient at high speeds; an electric motor was to be attached to the main reduction gear to turn the drive shaft and propel the ship at speeds under , such as during BMD or maritime security operations. Use of the HED for half the time could extend time on station by 2.5 days before refueling. In March 2018, the Navy announced the HED would be installed on to test the technology, but upgrades of further destroyers would be halted due to changed budget priorities.
Also in 2016, four destroyers of the U.S. 6th Fleet based in Naval Station Rota, Spain (USS Carney, USS Ross, USS Donald Cook, and USS Porter) received self-protection upgrades, replacing one of their two Phalanx CIWS with a SeaRAM CIWS, which combines the Phalanx sensor dome with an 11-cell RIM-116 launcher. This was the first time the system was paired with an Aegis ship. Another four ships (USS Arleigh Burke, USS Roosevelt, USS Bulkeley, and USS Paul Ignatius) have since been forward-deployed to Rota and also received a SeaRAM.
In February 2018, Lockheed Martin received a contract to deliver its High Energy Laser with Integrated Optical-dazzler and Surveillance (HELIOS) system for installation onto an Arleigh Burke destroyer. HELIOS is a "60+ kW"-class laser, scalable to 120 kW, that can "dazzle" or destroy small boats and UAVs up to away. It would be the first laser weapon put on a warship. In November 2019, had the Optical Dazzling Interdictor, Navy (ODIN) system installed. ODIN differs from the XN-1 LaWS previously mounted on in that ODIN functions as a dazzler, which blinds or destroys optical sensors on drones rather than shooting down the aircraft. HELIOS was delivered to the Navy in August 2022 and installed on . Preble was expected to begin at-sea testing of the HELIOS in FY2023.
Also by 2018, all Arleigh Burke-class ships homeported in the Western Pacific were scheduled to have upgraded ASW systems, including the TB-37U MFTA replacing the AN/SQR-19 TACTAS.
In FY2019, the Navy started a program to procure the Mod 4 variant of the Mark 38 machine gun system to address "unmanned aerial systems (UAS) and high speed maneuverable unmanned surface vehicle (USV) threats." Mod 4 will incorporate the 30 mm Mk44 Bushmaster II instead of the 25 mm M242 Bushmaster of previous variants. The Mk 38 Mod 4 is planned to be fielded on Flight IIA and III Arleigh Burke-class destroyers.
In October 2020, National Security Advisor Robert C. O'Brien said that all three Flights of the Arleigh Burke-class destroyer would field the Common-Hypersonic Glide Body (C-HGB) missile developed under the Conventional Prompt Strike program. However, the C-HGB is expected to be around wide, making it too large to fit in Mk 41 VLS tubes or on deck launchers. Installing them on Arleigh Burke destroyers would require removing some Mk 41 cells to accommodate the larger weapon, an expensive and time-consuming process. There is criticism of this idea: the oldest Flight I ships would need a service life extension to justify refit costs that would only prolong their service lives a short time when they are already more expensive to operate, and the newest Flight III ships that are optimized for BMD would be given a new, complex mission requiring a major refit shortly after introduction.
About 20 Flight IIA destroyers will undergo further modernization under the DDG MOD 2.0 program. DDG MOD 2.0 will backfit SPY-6(V)4 and Aegis Baseline 10 to provide similar capabilities to Flight III ships, as well as upgrade cooling systems to support the new radar. DDG MOD 2.0 will also deliver the AN/SLQ-32(V)7 EW suite, which adds the SEWIP Block 3 electronic attack subsystem. In May 2021, the Navy approved a "Smart Start Plan" for four ships—DDGs 91, 93, 95, 97—to make a gradual transition to DDG MOD 2.0. These ships will undergo a DDG MOD 1.5 phase that provides the SLQ-32(V)7; in 2023, DDG 91 became the first destroyer to receive SLQ-32(V)7. They will then receive the SPY-6(V)4, Aegis Baseline 10, and cooling system upgrades during a later depot modernization period.
Starting in 2025, the Navy will replace Phalanx CIWS on the destroyers with RIM-116 Rolling Airframe Missile (RAM) launchers to improve their point defense capability. Arleigh Burkes with the latest Aegis baselines will receive the 21-cell Mk 49 RAM launcher; Arleigh Burkes with older Aegis software will receive the 11-cell SeaRAM. It is unclear if ships with two Phalanx CIWS or ships already in a Phalanx-SeaRAM configuration will retain one Phalanx.
Production restarted
was originally intended to be the last of the Arleigh Burke class. The Navy planned to shift production to the Zumwalt-class destroyer focusing on NGFS and littoral operations. However, at a July 2008 hearing, Navy officials announced intentions to restart Arleigh Burke production in place of additional Zumwalts, testifying to the latter's inability to counter emerging ballistic missiles, anti-ship missiles, and blue-water submarines. Arleigh Burke-class destroyers have been in production for longer than any other surface combatant class in the U.S. Navy's history.
In April 2009, the Navy announced a plan limiting the Zumwalt class to three units while ordering another three Arleigh Burke-class ships from both Bath Iron Works and Ingalls Shipbuilding. In December 2009, Northrop Grumman received a $170.7 million letter contract for long lead-time materials. Shipbuilding contracts for DDG 113 to DDG 115 were awarded in mid-2011 for $679.6 million–$783.6 million; these do not include government-furnished equipment such as weapons and sensors, which took the average cost of the FY2011/12 ships to about $1.843 billion per vessel.
DDG 113-115 are "restart" ships, similar to previous Flight IIA ships, but including modernization features such as OACE and the TB-37U MFTA, which are being backfit onto previous ships.
The U.S. Navy was considering extending the acquisition of Arleigh Burke-class destroyers into the 2040s, according to revised procurement tables sent to Congress, with the procurement of Flight IV ships from 2032 through 2041. This was canceled to cover the cost of the s, with the air defense commander role retained on one cruiser per carrier strike group.
In April 2022, the Navy proposed a procurement plan for nine ships, with an option for a tenth, to build two ships a year from 2023 to 2027. Some lawmakers pushed to add a third ship to be built in 2023, bringing the total of the proposed deal to eleven ships. This would follow the Navy's two-ship per year procurement from 2018 to 2022.
Flight IIA Technology Insertion
DDG-116 to DDG-124 and DDG-127 will be "Technology Insertion" ships with elements of Flight III. For example, and onwards have the AN/SPQ-9B, a feature of Flight III, instead of the AN/SPS-67. Flight III proper began with the third ship procured in 2016, (DDG-125).
Flight III
In place of the canceled CG(X) program, the U.S. Navy began detailed design work on a DDG 51 Flight III design in FY2013. The Navy planned to procure 24 Flight III ships from FY2016 to FY2031. In June 2013, it awarded $6.2 billion in destroyer contracts. Costs for the Flight III ships increased as requirements for the program grew, particularly related to the planned Air and Missile Defense Radar (AMDR) needed for the IAMD role. An AMDR with a mid-diameter of had been proposed for CG(X), while the DDG 51 Flight III design could carry an AMDR with a mid-diameter of only . The Government Accountability Office (GAO) found that the design would be "at best marginally effective" because of the "now-shrunken radar". The U.S. Navy disagreed with the GAO findings, stating that the DDG 51 hull was "absolutely" capable of fitting a large enough radar to meet requirements.
The Flight III's AN/SPY-6 AMDR with a mid-diameter of uses an active electronically scanned array with digital beamforming, compared to the previous passive electronically scanned array AN/SPY-1D with a mid-diameter of . According to the SPY-6's contractor Raytheon, the 37-RMA SPY-6(V)1 is 30 times more sensitive and capable of detecting objects "half the size at twice the distance" compared to the SPY-1D. The Flight III's SPY-6 is integrated with Aegis Baseline 10. The new radar also requires more power; the three-megawatt, 450 V AG9140 generators were upgraded to four-megawatt, 4,160 V AG9160 generators. Additionally, the air conditioning plants were upgraded to increase the ships' cooling capacity. The area near where the two rigid-hull inflatable boats (RHIBs) are stored was enclosed to accommodate additional crew, so the RHIBs are stacked. Other modifications include replacement of the Halon-based fire suppression system with a water mist system and strengthening of the hull to support the design's additional weight.
14 Flight III ships have been ordered, and Flight III IOC is planned for 2024. The U.S. Navy may procure up to 42 Flight III ships for an overall total of 117 ships of the class.
Replacement
In April 2014, the U.S. Navy began the development of a new destroyer to replace the Arleigh Burke class called the "Future Surface Combatant". The new class is expected to enter service in the 2030s and initially serve alongside the Flight III Arleigh Burkes. The destroyer class will incorporate emerging technologies like lasers, onboard power-generation systems, increased automation, and next-generation weapons, sensors, and electronics. They will use technologies from other platforms, such as the Zumwalt-class destroyer, littoral combat ships, and the .
The Future Surface Combatant may place importance on the Zumwalt-class destroyer's electric drive system that provides propulsion while generating 58 megawatts of electrical power, levels required to operate future directed energy weapons. Initial requirements for the Future Surface Combatant emphasize lethality and survivability. The ships must also be modular to allow for inexpensive upgrades of weaponry, electronics, computing, and sensors over time as threats evolve. The Future Surface Combatant has evolved into the Large Surface Combatant, which became the DDG(X). The Navy plans to procure the first DDG(X) in FY2032.
Operational history
The class saw its first combat action through Tomahawk Land Attack Missile (TLAM) strikes against Iraq. Over 3 and 4 September 1996, and launched thirteen and eight TLAMs, respectively, as part of Operation Desert Strike. In December 1998, Arleigh Burke-class destroyers again performed TLAM strikes as part of Operation Desert Fox. Eleven Arleigh Burkes supported carrier strike groups engaged in Operation Iraqi Freedom, which included TLAM launches against ground targets during the operation's opening stages in 2003.
In October 2011, the Navy announced that four Arleigh Burke-class destroyers would be forward-deployed in Europe to support the NATO missile defense system. The ships, to be based at Naval Station Rota, Spain, were named in February 2012 as Ross, Donald Cook, Porter, and Carney. By reducing travel times to station, this forward deployment allows for six other destroyers to be shifted from the Atlantic in support of the Pivot to East Asia. Russia threatened to quit the New START treaty over this deployment, calling it a threat to their nuclear deterrent. In 2018, CNO Admiral John Richardson criticized the policy of keeping six highly mobile BMD platforms "in a little tiny box, defending land", a role that he believed could be performed equally well at less cost by shore-based systems.
In October 2016, the Arleigh Burke-class destroyers Mason and Nitze were deployed to the coast of Yemen after a UAE auxiliary ship was struck in an attack for which Houthi rebels claimed responsibility. On 9 October, while in the Red Sea, Mason detected two anti-ship missiles headed toward herself and nearby USS Ponce fired from Houthi-controlled territory. Mason launched two SM-2s, one ESSM, and a Nulka decoy. One AShM was confirmed to have struck the water on its own, and it is unknown if the second missile was intercepted or hit the water on its own. On 12 October, in the Bab el-Mandeb strait, Mason again detected an inbound anti-ship missile, which was intercepted at a range of by an SM-2. On 13 October, Nitze conducted TLAM strikes destroying three Houthi radar sites used in the previous attacks. Back in the Red Sea, Mason experienced a third attack on 15 October with five AShMs. She fired SM-2s and decoys, destroying or neutralizing four missiles. Nitze neutralized the fifth missile with a radar decoy.
On 7 April 2017, the Arleigh Burke-class destroyers Ross and Porter conducted a TLAM strike against Shayrat Airfield, Syria, in response to Syrian President Bashar Assad's chemical attack on his people three days prior. The ships fired a total of 59 Tomahawk missiles. On 14 April 2018, Laboon and Higgins conducted another TLAM strike against Syria. They fired seven and twenty-three TLAMs, respectively. The strike targeted chemical weapon sites as part of a continued effort against Assad's use of chemical warfare. The Arleigh Burke-class destroyers Donald Cook and Winston S. Churchill took positions in the Mediterranean prior to the 2018 strike to mislead defending forces.
In October and November 2023, the Arleigh Burke-class destroyers Carney and Thomas Hudner, while deployed in the Red Sea, shot down numerous drones and missiles. On 19 October, Carney shot down at least three cruise missiles and eight drones that were potentially targeting Israel. On 15 and 22 November, Thomas Hudner shot down numerous drones launched by Houthi rebels from Yemen. On 27 November, Carney detected two ballistic missile launches from Houthi-controlled territory headed towards herself and nearby M/V Central Park; they splashed ten nautical miles away. On 29 November, Carney intercepted another Houthi missile. On 30 December, Gravely shot down two anti-ship ballistic missiles fired from Houthi-controlled territory at herself and nearby container ship Maersk Hangzhou. On 30 January 2024, a Houthi anti-ship cruise missile fired toward the Red Sea came within one mile of Gravely; she used her Phalanx CIWS to shoot down the missile.
During the Iranian strikes on Israel on 13 April 2024, USS Arleigh Burke and USS Carney fired four to seven SM-3s, shooting down at least three Iranian ballistic missiles. This was the first time the SM-3 was employed in combat.
On 1 October 2024, USS Bulkeley and USS Cole fired 12 SM-3 and SM-6 missiles against Iranian ballistic missiles.
Accidents and major incidents
USS Cole bombing
was damaged on 12 October 2000 in Aden, Yemen, while docked by an attack in which a shaped charge of 200–300 kg in a boat was placed against the hull and detonated by suicide bombers, killing 17 crew members. The ship was repaired and returned to duty in 2001.
USS Porter and MV Otowasan collision
On 12 August 2012, USS Porter collided with the oil tanker MV Otowasan near the Strait of Hormuz; there were no injuries. The U.S. Navy removed Porters commanding officer from duty. Repairs took two months at a cost of $700,000.
USS Fitzgerald and MV ACX Crystal collision
On 17 June 2017, collided with the MV ACX Crystal cargo ship near Yokosuka, Japan. Seven sailors drowned. Following an investigation, the ship's commanding officer, executive officer, and Command Master Chief Petty Officer were relieved of their duties. In addition, close to a dozen sailors were given non-judicial punishment for losing situational awareness. Repairs were originally to be completed by the summer of 2019. However, initial repairs were made by February 2020. After the subsequent sea trials, she was brought in for additional repairs. The ship departed for her home port in June 2020.
USS John S. McCain and Alnic MC collision
On 21 August 2017, USS John S. McCain collided with the container ship Alnic MC. The collision injured 48 sailors and killed 10, whose bodies were all recovered by 27 August. The cause of the collision was determined to be poor communication between the two ships and the bridge crew lacking situational awareness. In the aftermath, the ship's top leadership, including the commanding officer, executive officer, and Command Master Chief Petty Officer, were removed from command. In addition, top leadership of the U.S. Seventh Fleet, including the commander, Vice Admiral Joseph Aucoin, were relieved of their duties due to a loss of confidence in their ability to command. Other commanders who were relieved included Rear Admiral Charles Williams, commander of Task Force 70, and Captain Jeffrey Bennett, commodore of Destroyer Squadron 15. This was the third incident involving a U.S. Navy ship in 2017, with a repair cost of over $100 million.
Contractors
Builders: 38 units constructed by General Dynamics, Bath Iron Works Division, and 35 by Huntington Ingalls Industries (formerly Northrop Grumman Ship Systems), Ingalls Shipbuilding
AN/SPY-1 radar and Aegis Combat System integrator: Lockheed Martin
AN/SPY-6 radar: Raytheon
Ships in class
Derivatives
Destroyer classes based on the Arleigh Burke have been adopted by the following naval forces:
The Japanese Maritime Self-Defense Force:
The Republic of Korea Navy:
In popular culture
The 2009 film Transformers: Revenge of the Fallen features USS Preble.
The 2012 film Battleship features , , and .
The 2013 film Captain Phillips features USS Truxtun, which stood in for the ship from the true event, USS Bainbridge.
The 2014 television series The Last Ship, loosely based on the 1988 novel of the same name, is set on the fictional . Its hull designation in the book is DDG 80, but it was changed to DDG 151 for the television series to avoid confusion with the real-life USS Roosevelt, which did not exist when the book was written. , a Flight IIA Arleigh Burke-class destroyer, stood in for Nathan James during filming.
The 2015 film San Andreas features an unidentified Arleigh Burke-class destroyer during a scene following a tsunami. A second and third unidentified destroyer also appear near the end of the film in the San Francisco Bay.
| Technology | Naval warfare | null |
197407 | https://en.wikipedia.org/wiki/Red-winged%20blackbird | Red-winged blackbird | The red-winged blackbird (Agelaius phoeniceus) is a passerine bird of the family Icteridae found in most of North America and much of Central America. It breeds from Alaska and Newfoundland south to Florida, the Gulf of Mexico, Mexico, and Guatemala, with isolated populations in western El Salvador, northwestern Honduras, and northwestern Costa Rica. It may winter as far north as Pennsylvania and British Columbia, but northern populations are generally migratory, moving south to Mexico and the Southern United States. Claims have been made that it is the most abundant living land bird in North America, as bird-counting censuses of wintering red-winged blackbirds sometimes show that loose flocks can number in excess of a million birds per flock and the full number of breeding pairs across North and Central America may exceed 250 million in peak years. It also ranks among the best-studied wild bird species in the world. The red-winged blackbird is sexually dimorphic; the male is all black with a red shoulder and yellow wing bar, while the female is a nondescript dark brown. Seeds and insects make up the bulk of the red-winged blackbird's diet.
Taxonomy
The red-winged blackbird is one of five species in the genus Agelaius and is included in the family Icteridae, which is made up of passerine birds found in North and South America. The red-winged blackbird was formally described as Oriolus phoeniceus by Carl Linnaeus in 1766 in the twelfth edition of his Systema Naturae, but was later moved with the other American blackbirds to the genus Agelaius by Louis Pierre Vieillot in 1816. Linnaeus specified the type location as "America" but this was restricted to Charleston, South Carolina in 1928. The genus name is derived from Ancient Greek , meaning "gregarious". The specific epithet, , is Latin meaning "crimson" or "red". The red-winged blackbird is a sister species to the red-shouldered blackbird (Agelaius assimilis) that is endemic to Cuba. These two species are together sister to the tricolored blackbird (Agelaius tricolor) that is found on the Pacific coast region of the California and upper Baja California in Mexico.
Depending on the authority, between 20 and 24 subspecies are recognised which are mostly quite similar in appearance. However, there are two isolated populations of bicolored blackbirds that are quite distinctive: A. p. californicus of California and A. p. gubernator of central Mexico. The taxonomy and relationships between these two populations and with red-winged blackbirds is still unclear. Despite the similarities in most forms of the red-winged blackbird, in the subspecies of Mexican Plateau, A. p. gubernator, the female's veining is greatly reduced and restricted to the throat; the rest of the plumage is very dark brown, and also in a different family from the European redwing and the Old World common blackbird, which are thrushes (Turdidae). In the California subspecies, A. p. californicus and A. p. mailliardorum, the veining of the female specimens also covers a smaller surface and the plumage is dark brown, although not in the gubernator grade; and also its superciliary list is absent or poorly developed. The male subspecies mailliardorum, californicus, aciculatus, neutralis, and gubernator lack the yellow band on the wing that is present in most male members of the other subspecies. The red-shouldered blackbird was formerly considered as a subspecies of red-winged blackbird. They were split by the American Ornithologists' Union in 1997.
Listed below are the red-winged blackbird subspecies and groups of subspecies recognized as of January 2014 with their respective distribution areas and the location of their wintering quarters:
Group phoeniceus:
A. p. arctolegus – from southeastern Alaska and Yukon to the north-central United States; migrates to the south-central United States.
A. p. fortis – from Montana to southeastern New Mexico (east of the Rocky Mountains); migrates to Texas.
A. p. nevadensis – from southeastern British Columbia to Idaho, southeastern California, and southern Nevada; migrates to southern Arizona.
A. p. caurinus – from the southwest coast of British Columbia to northwest California; migrates to central California.
A. p. aciculatus – mountains of central-southern California (central-eastern Kern County).
A. p. neutralis – from the coast of southern California (San Luis Obispo County) to the northwest of Baja California.
A. p. sonoriensis – from southeastern California to northeast Baja California, southern Nevada, central Arizona, and northwest Mexico.
A. p. nyaritensis – coastal plains of southwestern Mexico (Nayarit).
A. p. grinnelli – Pacific Slope from western Guatemala to northwest Costa Rica (Guanacaste).
A. p. phoeniceus – from southeastern Canada to Texas and the southeastern United States.
A. p. littoralis – Gulf Coast from Southeast Texas to Northwest Florida.
A. p. mearnsi – extreme southeast Georgia and northern Florida.
A. p. floridanus – South Florida (Everglades to Key West).
A. p. megapotamus – from central Texas and the lower Rio Bravo valley to eastern Mexico (northern Veracruz).
A. p. richmondi – Caribbean slope from Mexico (southern Veracruz) to Belize and northern Guatemala.
A. p. pallidulus – southeast Mexico (north of the Yucatán Peninsula).
A. p. nelsoni – south-central Mexico (from Morelos and the adjacent Guerrero west of Puebla and Chiapas).
A. p. matudae – tropical southeast Mexico (not always recognised).
A. p. arthuralleni – northern Guatemala.
A. p. brevirostris – Caribbean slope of Honduras and southeastern Nicaragua (sometimes included in richmondi).
A. p. bryanti – northwest Bahamas.
Group californicus/ mailliardorum:
A. p. mailliardorum – Central California Coast.
A. p. californicus – Central Valley of California.
Group gubernator :
A. p. gubernator – Mexican Plateau (Durango to Zacatecas, Mexico and Tlaxcala).
Description
The common name for the red-winged blackbird is taken from the mainly black adult male's distinctive red shoulder patches, or epaulets, which are visible when the bird is flying or displaying. At rest, the male also shows a pale yellow wingbar. The spots of males less than one year old, generally subordinate, are smaller and more orange than those of adults. The female is blackish-brown and paler below. The female is smaller than the male, at long and weighing , against his length of and weight of . The smallest females may weigh as little as whereas the largest males can weigh up to . Each wing can range from , the tail measures , the culmen measures and the tarsus measures . The upper parts of the female are brown, while the lower parts are covered by an intense white and dark veining; also presents a whitish superciliary list. According to Crawford (1977), females exhibit a salmon-pink stain on the shoulders and a light pink color on the face and below this when they are a year old, while older females show a stain usually more crimson on the shoulders and a darker pink hue on and under the face. Observations in females in captivity indicate that small amounts of yellow pigment are present on the shoulders of these after leaving the nest, that the concentration of the pigment increases with the first winter plumage after the change of the feathers and that the passage from yellow to orange generally takes place in the second summer with the acquisition of the second winter plumage, after which no further changes in feather color occur. The colored area on the wing increases in surface with the age of the female, and varies in intensity from brown to a bright red-orange similar to that of the males in their first year.
Young birds resemble the female, but are paler below and have buff feather fringes. Both sexes have a sharply pointed bill. The tail is of medium length and is rounded. The eyes, bill, and feet are all black. Unlike most North American passerines which develop their adult plumage in their first year of life so that the one-year-old and the oldest individual are indistinguishable in the breeding season, the red-winged blackbird does not. It acquires its adult plumage only after the breeding season of the year following its birth when it is between thirteen and fifteen months of age. Young males go through a transition stage in which the wing spots have an orange coloration before acquiring the most intense tone typical of adults.
The male measures between in length, while the female measures . Its wingspan is between approximately. Both the peak male and the legs, the claws and the eyes are black; in the female beak is dark brown and clear in the upper half at the bottom, and the tail is medium in length and rounded. As in other species polygyny exists, the red-winged blackbird exhibit considerable sexual dimorphism both in plumage and size, males weighing between 65 and 80g the females about 35g. Males are 50% heavier than females, 20% larger in its linear dimensions, and 20% larger compared to the length of their wings. The trend towards greater dimorphism in the size of non-monogamous ichterid species indicates that the larger size of males has evolved due to sexual selection.
The male is unmistakable except in the far west of the US, where the tricolored blackbird occurs. Males of that species have a darker red epaulet edged with white, not yellow. Females of tricolored, bicolored, red-shouldered and red-winged blackbirds can be difficult to identify in areas where more than one form occurs. In flight, when the field marks are not easily seen, the red-winged blackbird can be distinguished from less closely related icterids such as common grackle and brown-headed cowbird by its different silhouette and undulating flight.
Staining of red spots and yellow bands
Two keto-carotenoids – carotenoid with a ketone group – reds synthesized by the birds themselves – namely astaxanthin and canthaxanthin – are responsible for the bright red color of the wing spots, but two yellow dietary precursor pigments – lutein and zeaxanthin – are also present in moderately high concentrations in red feathers. Astaxanthin is the carotenoid more abundant (35% of the total), followed by lutein (28%), canthaxanthin (23%) and zeaxanthin (12%). Such a balanced combination of dietary precursors and metabolic derivatives in colored feathers is quite unusual, not only within a species as a whole but also in individuals and particular feathers.
After removing the carotenoids in an experiment, the red feathers acquired a deep brown coloration. This is because the feather barbles of the colored spot contain melanin pigments – mainly eumelanin, which was equivalent to 83% of all melanins, but also pheomelanin at a concentration approximately equal to that of carotenoids, which seems to be a rare trait for carotenoid-based ornamental plumage. On the other hand, the feathers of the yellow stripes of the males are devoid of carotenoids —except occasionally when they appear tinged with a pink coloration derived from small amounts of said pigments— and present high concentrations of pheomelanin —82% of all melanins. Melanin concentrations in the yellow band are even higher than in the red spot.
Role of wing spots
These markings are vital in the defense of the territory. Males with larger spots are more effective at chasing away their non-territorial rivals and are more successful in contests within aviaries. When the red shoulder patches were dyed black as part of an experiment, 64% of males lost their territories, while only 8% of control subjects did. However, males whose wings had been dyed before they had mated could still attract females and successfully reproduce. In the red-winged blackbird, the spots on the wings are a sign of threat among males and have an unimportant role, if any, in intersex encounters. Therefore, the spots are likely to have evolved in response to pressures linked to intrasexual selection. Additionally, neither the size nor the coloration of the same are linked to the reproductive success of the males with those females that are not their mates, that is, those with which they eventually mate. It has been suggested that also in the case of females, the best explanation for the evolution of a variable coloration in the shoulder spots is that their intensity indicates their physical condition in aggressive encounters between them.
The fact that female red-winged blackbirds do not appear to consistently use variability in the size and color of male wing spots when choosing a mate runs counter to the classic role of carotenoid pigmentation ornamental feathers, mostly used in the attraction of a couple. In turn, its use as a sign of aggressiveness and social status against rival males is not a common trait in carotenoid ornaments. On the other hand, ornaments with a preponderance of melanin do tend to have an important role as an indicator of status in avian populations, so the spots of red-winged blackbirds seem to work more as melanin ornaments than as carotenoids.
Vocalizations
The calls of the red-winged blackbird are a throaty check and a high slurred whistle, . The male's song, accompanied by a display of his red shoulder patches, is a scratchy , except that in many western birds, including bicolored blackbirds, it is . The female also sings, typically a scolding chatter .
Feather molt
The most critical period of feather molting runs from late August to early September. When viewed in flight, they have a misaligned or "moth-eaten" appearance and generally slower and more laborious travel. Their mobility is reduced due to the lack of several remiges or rectrices or these are not entirely renewed. Most of the red-winged blackbirds have moved almost entirely by October. By then, some birds have not completed the molt of the feathers of the capital region and the helmsmen of the center of the tail and the internal secondary sprouts have only partially emerged from the pod. Virtually all individuals have completed their molts by mid-October.
Birds do not begin their migration to wintering quarters until the two outer primary sprouts and the two inner or central rectrices have completed at least two-thirds of their development. Therefore, there is a correlation between molting, particularly replacement of the remiges and rectrices, and fall migration in red-winged blackbirds.
Succession of plumages and molts
Juvenile plumage: in both males and females, juvenile plumage is similar in color and pattern to that of the adult female, except for the yellow coloration of the lower body and the sides of the head . The shoulder spots appear mottled and are brown and yellow or brown and beige. Other sources describe the juvenile plumage of the male as follows: above, including the sides of the head, the wings, the tail and the lesser coverts – the feathers of the "shoulders" -, dull brownish black – not seen no red at this stage—, with beige-bordered feathers, paler and narrower in primaries, rectrices, head, and bladderand deeper in the scapulars and secondaries; below, pinkish beige, ocher above chin, densely veined — except chin — with brownish black stripes; undefined ocher-beige superciliary list. The female is described as more brown above with less beige below and a closely veined chin.
First pre-basic molt (or post-juvenile molt): Generally begins 45–60 days after individuals have left the nest. The molt begins in August and the time at which it begins varies between early and late clutches. This is a complete molt, except for the retention in some individuals of a few down feathers under the wing.
Appearance in autumn and winter of the first basic plumage (also called the first winter plumage or immature plumage): in the male, it is black, with the feathers of the upper part bordered with brown or beige and those of the lower part bordered with beige or white. The bird then has a mottled appearance. The shoulder spot is generally orange with black speckles, especially on the yellow stripe. Some immature males have a reddish patch like adults, but with black specks on the yellow band. Others have a blackish spot on the shoulder. Other sources describe the male's first winter plumage as follows: all plumage, including wings and tail, greenish-black, much of it with beige and rust edges, paler underneath and faint or absent in primaries and rectrices; minor wing covers – "shoulders" – orpiment-dull orange, each feather with black subterminal bars or spots; medium beige, deep ocher coverts, generally mottled with black in subterminal areas, mainly in the beards of the inner part of the feathers; generally black spines.
Females are dark above with feathers bordered with beige or rust. The lower part shows black and white streaks, but with more beige on the chest and sometimes on the flanks than in the reproductive plumage. The medium and secondary major cover feathers are noticeably bordered with beige. They generally lack the pink coloration on the chin and throat – which may be beige, yellow, or light salmon – and the crimson wing patch – which may be rust, orange, or gray – typical of second winter plumage.
Appearance in spring and summer of the first basic plumage (also known as the first breeding plumage or sub-adult plumage): the males, which normally are not yet defending a territory, present a dull black coloration acquired by the wear of the edges of the feathers. The stain alar may be more noticeable than in the first winter plumage due to wear subterminal blacklists, which generally only remain in the form of small black spots. Mottled wing spots are characteristic of young individuals and the extent of orange is highly variable. Females, which are usually already breeding, exhibit plumage similar to the first winter plumage, but darker above due to wear on beige edges or rust on feathers. The chest is less beige, and has a black and white veining.
Second and subsequent prebasic molts (or postnuptial molts) : the second prebasic molt takes place approximately one year after the first. Postnuptial molts constitute complete molts, except for the occasional retention of a few down feathers under the wing.
Appearance in autumn and winter of the second and subsequent basic plumages (also called adult winter plumage or second winter plumage) : the male no longer exhibits the mottled appearance characteristic of the autumn and winter appearance of the first basic plumage. The lower part is almost immaculate and similar to that of the reproductive plumage. The head and back feathers and the secondary covers are bordered with brown and beige. The wing blot becomes bright scarlet-vermilion, while deep ocher beige coloration appears on the medium coats. Other sources describe the adult winter plumage of the male as follows: glossy greenish black; feathers of the head, back, major wing covers and tertiary more or less bordered, depending on the individual, beige and rust; below, paler or absent edges. The female's plumage is similar to the autumn and winter appearance of the first basic plumage, except that the wing spot is generally crimson and the chin and throat pink.
Aspect in summer and spring of the second and subsequent basic plumages (or adult bridal plumage): it is acquired through the wear of the feathers. is similar to the appearance in the spring and summer of the first basic plumage, but the alar spots of both sexes have a more intense color and chin and throat of the female. The beige and brown edges of the male feathers disappear.
Despite the fact that the brown or white tips on the feathers of the males are larger just after the molt and that they wear out throughout the year, the individuals vary considerably regarding the size of these non-black tips on the feathers in spring.
Wing feathers
Complete replacement of wing feathers takes about eight weeks. However, birds in their first year frequently retain some of the under-wing coverts and juvenile tertiary remiges after post-juvenile moulting. Of seventy immature males examined during the last week of October, it retained 70% of older lower primary blankets. In most cases where partial replacement of the covert feathers occurs, it is the proximal coverts that the bird retains.
Remiges
Primary remiges are one of the first feathers to molt. The molt of these feathers proceeds regularly from the innermost primary—primary I—to the outermost—primary IX. By October 1, most birds have either acquired the three new external primaries—VII, VIII, and IX—or they are at some advanced stage of development. The average dates for completion of the development of the new primary springs are: August 15 for primary I; September 1, primary II-IV; September 15, primary V and VI; and October 1, primaries VII-IX.
The molting of the secondary remiges begins with the most external—secondary I— and proceeds inward to secondary VI. The secondary I sheath appears around the same time that all the secondary covers have been replaced and rarely before mid-August. These feathers are not completely renewed until the beginning of October.
The molt of the tertiary remiges begins more or less at the same time as that of the secondary ones. The middle tertiary falls first, followed by the internal tertiary. Both feathers are often well developed again before the outer tertiary leaves the sheath.
Cover feathers and alula
The major primary blankets are changed along with their respective primary springs. Unlike the major primary coats, the major secondary coats molt earlier than the secondary spruces. The molting of these feathers is rapid, with several of them at the same stage of development simultaneously. The progression of the molt in these feathers is from the outside to the inside, as in the secondary remiges. Most birds has completed the change of the secondary coverts to August 15, more or less at the time is appreciable only secondary sheath remige I.
The molting of the lesser coverts begins early, often being the first feathers to fall. The onset of molting in male juveniles is particularly noticeable because it involves the replacement of the minor covers and results in the appearance of the reddish or orange wing spot. The new wing spot contrasts sharply with the yellowish-brown juvenile plumage in this area of the wing. The move of the minor blankets has generally been completed by September 1.
Alula feathers complete their development at about the same time as the last three primary sprouts. The marginal covers on the upper or outer surface of the forearm, located below the alula, shed at approximately the same time that the primary remix VI is being replaced.
The first feathers under the wing to molt are the marginal coverts, under the forearm. The shedding of these feathers begins at about the same time that the primary remix IV falls and is followed by that of the lower middle primary and lower middle secondary coats. The progression of the molting of the lower middle secondary blankets is from the outside to the inside, while that of the lower middle primary mats seems to be irregular or almost simultaneous. The medial lower coverts molt before the primary remiges VIII and IX. The major lower primary coverts and major lower secondary coverts molt last. The progression of the molting of these last feathers is the same as in the primary and secondary sprouts, that is, from the inside out and from the outside in, respectively.
Flow and capital pens
The caudal feathers comprise the rudder or rectrix feathers and the upper and lower tail covers. The tail covers begin to shed before the rectrices. Generally, the upper tail covers begin to shed first. Certain birds lose some rectrices by the end of the third week of August. The helmsmen in the center of the queue are the last rectrices to be renovated.
Molting in the capital region involves changing the feathers of the pileus and the sides of the head. It is one of the last parts of the body to begin feather replacement, but the renewal of most of the capital feathers is complete before that of the secondary feathers, tail feathers, and under-wing feathers. The beginning of the molt in this region coincides with the beginning of the development of the primary remige V or VI. Some individuals have already started replacing the capital feathers by mid-August. The molt begins at the pileus and the last areas of the capital region to complete it are the eye strip and the cheeks (malar region).
Other pens
In some birds, the first signs of molt in the ventral feathers appear during the last days of July, when the feathers of the anterior portion of the laterals begin to fall. From there, the molt progresses backward along the sides and forward toward the throat and chin. The last ventral feathers to be replaced are those towards the center of the abdomen. The molt of the dorsal feathers begins around the first week of August. It begins at the bladder, progresses to the upper back, and then to the cervical region.
The earliest evidence of molting in the humeral plumes corresponds to the last days of July. The molt comes from the anterior region backwards. The change of the femoral feathers begins later than that of the humeral ones. However, the progression is similar. Replacement of feather feathers rarely begins before August 15. The progression is generally from the proximal end of the tibia to the tarsometatarsal region.
Distribution and habitat
The red-winged blackbird is widely spread throughout North America, except in the arid desert, high mountain ranges, and arctic or dense afforestation regions. It breeds from central-eastern Alaska and Yukon in the northwest, and Newfoundland in the northeast, to northern Costa Rica in the south, and from the Atlantic to the Pacific. Northern populations migrate to the southern United States, but those that breed there, in Mexico, and in Central America are sedentary. Red-winged blackbirds in the northern reaches of the range are migratory, spending winters in the southern United States and Central America. Migration begins in September or October, but occasionally as early as August. In western and Central America, populations are generally non-migratory.
The red-winged blackbird inhabits open grassy areas. It generally prefers wetlands, and inhabits both freshwater and saltwater marshes, particularly if cattail is present. It is also found in dry upland areas, where it inhabits meadows, prairies, and old fields. In a large part of its distribution area, it constitutes the most abundant passerine bird in the swamps in which it nests. It is also present in areas without much water, where it inhabits open fields – often agricultural areas – and sparse deciduous forests.
In the winter of 1975–1976, near Milan in western Tennessee, red-winged blackbirds were observed resting in a mixed roost that came to house 11 million individuals in January and early February in a plantation of 4.5 hectares of yellow pine (Pinus taeda) with little undergrowth was seen in soybean fields during the day, being that these constituted only 21% of the habitat in the area and that the other bird species present in the roost were not commonly observed in these fields; they were also common in cornfields. the presences of the bird in feedlots has increased as winter progressed, but, accounting for less than 5% of the icterids and starlings recorded in both feedlots of cows and pigs, they were much rarer there than the brown-headed cowbird, common grackle and common starling.
Preferences by habitat types during the breeding season
During the breeding season, the density of breeding adults is much higher in swamps than in highland fields. Although the highest concentrations of nesting red-winged blackbird are found in swamps, most of them nest in upland habitats since they are much more abundant. In a study conducted in the Wood County, Ohio between 1964 and 1968, the density of territorial males in wetland habitat was found to be 2.89 times those in upland habitat. But, because of the small amount of wetland habitat, the total estimated upland population of territorial males
was 2.14 times the wetland population. Alfalfa (Medicago sativa) and other legume crops (hay) were the principal habitat for breeding redwings in the county. Despite the marked preference for wetland habitats, the greater population in uplands reflects wetlands' scarcity.
Vegetation type preferences during the breeding season
In that same study, reproductive individuals in upland habitats demonstrated a slight preference for old and new grasses such as Phleum pratense, Dactylis glomerata, Poa spp., Festuca spp., and Bromus spp. in the early part of the breeding season and new non-graminoid herbaceous plants in the mid and late season. In wetlands, they consistently preferred old and new broadleaf monocotyledons, primarily Carex spp, broadleaf, and Typha spp, and consistently rejected old and new narrow-leaved monocots, primarily narrow-leaved Phalaris arundinacea and Calamagrostis canadensis, and non-graminoid herbaceous plants. Red-winged blackbirds breeding in the highlands favored older tall vegetation only in the early part of the breeding season on April to May, and higher new vegetation and dense vegetation throughout the breeding season. For their part, those who settled in wetlands seemed to have a slight predilection for taller old vegetation. During the breeding season, this species is attracted to tall vegetation that restricts visibility.
Early season preferences for old grassland in the highlands and old broadleaf monocotyledons in wetlands point to the importance of upright residual vegetation. Highland grasses and broadleaf monocots in wetlands stand partially vertically and are easily visible in early spring, unlike clovers, narrow-leaved monocotyledons in wetlands and most non-graminoid herbaceous plants. Old alfalfa plants are also partially upright in early spring, but they are not as consistently chosen a species as old grasses. The start of the territorial activity is earlier when the amount of residual vegetation is large. Vegetation structural strength also appears to be important for nesting as females tend to broadleaf monocotyledons in wetlands throughout the breeding season and new non-graminoid herbaceous plants in the mid to late season.
In southwestern Michigan, the density of cattail stems in different ponds was positively related to the concentration of reproductive adults. However, other studies detected a predilection for more scattered vegetation. One found that red-winged blackbirds avoided the marsh wren (Pantaneros chivirenes), a predator common in nests of this species, reproducing among more dispersed vegetation, which was more easily defended from the chivirines; conversely, the chivirines seemed to prefer the denser vegetation, where they were more likely to avoid the aggressiveness of the red-winged blackbird; These differences in habitat selection between one species and another resulted in spatial segregation of their breeding areas. On the other hand, red-winged blackbirds tend to opt for small lots of vegetation and plants with thick stems.
Behavior
The red-winged blackbird is territorial, polygynous, gregarious and a short-distance migratory bird. Its way of flying is characteristic, with rapid wing flaps punctuated by brief periods of gliding flight. The behavior of males makes their presence easily perceived: they perch in high places such as trees, bushes, fences, telephone lines, etc. Females tend to stay low, prowling through the vegetation and building their nests. They can be found in home gardens, particularly during their migration, if seeds have been scattered on the ground. The forest curtains serve as a resting place during the day.
For several weeks after their first appearance in early spring, red-winged blackbirds are generally seen in flocks made up entirely of males. During those days, they are seldom seen at their breeding sites, except early morning and late afternoon. In most of the remaining hours of the day, they frequent open and often elevated agricultural land, where they feed mainly on grain stubble and grassy fields. When disturbed while eating, they fly to the nearest deciduous trees and immediately after landing they begin to sing.
Food
The red-winged blackbird is omnivorous. It feeds primarily on plant materials, including seeds from weeds and waste grain such as corn and rice, but about a quarter of its diet consists of insects and other small animals, and considerably more so during breeding season. It prefers insects, such as dragonflies, damselflies, butterflies, moths, and flies, but also consumes snails, frogs, eggs, carrion, worms, spiders and mollusks. The red-winged blackbird forages for insects by picking them from plants, or by catching them in flight. Sometimes insects are obtained by exploring the bases of aquatic plants with their small beaks, opening holes to reach the insects hidden inside. Aquatic insects, particularly emerging odonates, are of great importance in the diet of red-winged blackbirds that breed in swamps. These birds typically capture the odonates when the larvae climb up the stem of a plant from the water, get rid of their exuviae, and cling to the vegetation while their exoskeletons harden. The years of emergence of periodical cicadas provide an overabundant amount of food. According to Edward Howe Forbush, when red-winged blackbirds arrive north in the spring, they feed in the fields and meadows. Then, they follow the plows, collecting larvae, earthworms and caterpillars left exposed, and in case there is a plague of Paleacrita vernata caterpillars in a fruit orchard, these birds will fly a kilometer to get them for their chicks.
In season, red-winged blackbirds eat blueberries, blackberries, and other fruits. These birds can be lured to backyard bird feeders by bread and seed mixtures and suet. In late summer and in autumn, the red-winged blackbird will feed in open fields, mixed with grackles, cowbirds, and starlings in flocks which can number in the thousands. It feeds on corn while it is maturing; once the grain has hardened it is relatively safe from this bird, since its beak and digestive system are not adapted for the consumption of hard and whole corn grains, unlike the Common grackle, which has a longer and stronger beak. Studies of the stomachs of individuals of both sexes reveal that males consume higher proportions of crop grains, while females ingest a relatively larger amount of herb seeds and animal matter. In the winter of 1975–1976, near Milan, western Tennessee, corn and herb seeds were the main foods consumed by red-winged blackbird. Herbs whose seeds were commonly consumed were Sorghum halepense, Xanthium strumarium, Digitaria ischaemum, Sporobolus spp., Polygonum spp. and Amaranthus spp.
Breeding
The red-winged blackbird nests in loose colonies. The nest is built in cattails, rushes, grasses, sedge, or alder or willow bushes. The nest is constructed entirely by the female over the course of three to six days. It is a basket of grasses, sedge, and mosses, lined with mud, and bound to surrounding grasses or branches. It is located to above water.
A clutch consists of three or four, rarely five, eggs. Eggs are oval, smooth and slightly glossy, and measure . They are pale bluish green, marked with brown, purple, and/or black, with most markings around the larger end of the egg. These are incubated by the female alone, and hatch in 11 to 12 days. Red-winged blackbirds are hatched blind and naked, but are ready to leave the nest 11 to 14 days after hatching.
Red-winged blackbirds are polygynous, with territorial males defending up to 10 females. However, females frequently copulate with males other than their social mate and often lay clutches of mixed paternity. Pairs raise two or three clutches per season, in a new nest for each clutch.
The reproductive season of the red-winged blackbird extends approximately from the end of April to the end of July. On the other hand, in different states has been estimated that the period in which the active nests contained eggs lay between beginning in late April and early late August; and in northern Louisiana nests were found to harbor chicks from late April to late July. The peak of the nesting season (the time with the highest number of active nests) has been recorded between the first half of May and the beginning of June in different places. A study in eastern Ontario found that although red-winged blackbirds began nesting earlier in years with warm springs, associated with low winter values in the North Atlantic Oscillation Index, egg laying dates remained unchanged. Male testosterone levels peak in the early part of the breeding season, but remain high throughout the season. Females reproduce for up to ten years.
Many aspects of territorialism (for example, the number of male songs and displays, and the number of intrusions into foreign territories) peak before mating. After this, the frequency of many of the territorial behaviors decreases and the males are mainly concerned with defending the females, eggs, and chicks against predation. Experiments in the systematic removal of birds from their territories suggest that the extra population of males that is present in swamps before copulations disappears after copulation.
Predation of eggs and nestlings is quite common. Nest predators include snakes, mink, raccoons, and other birds, even as small as marsh wrens. The red-winged blackbird is occasionally a victim of brood parasites, particularly brown-headed cowbirds. Since nest predation is common, several adaptations have evolved in this species. Group nesting is one such trait which reduces the risk of individual predation by increasing the number of alert parents. Nesting over water reduces the likelihood of predation, as do alarm calls. Nests, in particular, offer a strategic advantage over predators in that they are often well concealed in thick, waterside reeds and positioned at a height of one to two meters. Males often act as sentinels, employing a variety of calls to denote the kind and severity of danger. Mobbing, especially by males, is also used to scare off unwanted predators, although mobbing often targets large animals and man-made devices by mistake. The brownish coloration of the female may also serve as an anti-predator trait in that it may provide camouflage for her and her nest while she is incubating.
Predators and parasites
Predators of red-winged blackbirds include such species as raccoons, American mink, long-tailed weasels, black-billed magpies, common grackles, owls, red-tailed hawks, short-tailed hawks, and snakes (such as the northern water snake and the plains garter snake). Ravens and grazers such as marsh wrens feed on eggs (and even small chicks), if the nest is left unattended, destroying the eggs, occasionally drinking from them, and pecking the nestlings to death.
The relative importance of different nest predators varies by geographic region: the top predators in different regions include the marsh wren in British Columbia, the magpies in Washington, and the raccoons in Ontario. The incidence of avian predation in red-winged blackbird nests is higher in western populations than in eastern populations.
Due to high predation rates, especially of eggs and chicks, the red-winged blackbird has developed various adaptations to protect its nests. One of them consists of nesting in groups, which reduces the danger since there is a greater number of alert parents. Nesting over water also lowers the chances of an attack. Nests in particular offer a strategic advantage as they are often hidden among dense riparian reeds, at a height of one or two meters. males often act as sentinels, using a repertoire of calls. Males in particular hunt down potential predators to scare them away, even when dealing with much larger animals. Aggressiveness of the red-winged blackbird towards the marsh wren, which also nests in swamps, causes a partial interspecies territorialism. On the other hand, nocturnal predators such as raccoons and American mink are not attacked by adults. Coloration of the female could serve to camouflage it, protecting it and its nest when it is incubated.
The red-winged blackbird can accommodate ectoparasites such as various Phthiraptera, the Ischnocera Philopterus agelaii and Brueelia ornatissima, and hematophagous mites, like Ornithonyssus sylviarum, and endoparasites like Haemoproteus quiscalus, Leucocytozoon icteris, Plasmodium vaughani, nematodes, flukes and tapeworms.
Territorial
The red-winged blackbird aggressively defends its territory from other animals, and will attack much larger birds. During breeding season, males will swoop at humans who encroach upon their nesting territory. Male red-winged blackbirds also exhibit important territorial behaviors, most of which provides them with the necessary fidelity for many years to come. A few important factors for male red-winged blackbirds’ adherence to territories include food, hiding spaces from predators, types of neighbors, and reactions towards predators. Additionally, a study was done on site fidelity and movement patterns by Les D. Beletsky and Gordon H. Orians in 1987 which explained much of the males’ territorial behaviors once migrated and settled onto a territory of their own. Sufficient evidence had shown that males are committed to staying in their territory over a long period of time and are not more likely to change territories at a younger age due to limited experience of knowledge for success. Studies also showed that most of the males that were first-time movers to a new territory were between two and three years old. The majority of males that moved were young and inexperienced. Later on they had moved towards more available territories. If males had chosen to leave their territory for reproductive success, as an example, they would do so within a short distance. Males who moved shorter distances were more successful in reproducing than those who moved longer distances. Further studies showed that when males moved further away from their territories there was a decrease in probability of successfully fledging.
The maximum longevity of the red-winged blackbird in the wild is 15.8 years.
Migration
Red-winged blackbirds that breed in the northern part of their range, i.e., Canada and border states in the United States, migrate south for the winter. However, populations near the Pacific and Gulf coasts of North America and those of Middle America are year-round resident. Red-winged blackbirds live in both Northern U.S. and Canada, ranging from Yucatan Peninsula in the south to the southern part of Alaska. These extensions account for the majority of the continent stretching from California's Pacific coast and Canada to the eastern seaboard. Much of the populations within Middle America are non-migratory. During the fall, populations begin migrating towards Southern U.S. Movement of red-winged blackbirds can begin as early as August through October. Spring migration begins anywhere between mid-February to mid-May. Numerous birds from northern parts of the U.S., particularly the Great lakes, migrate nearly between their breeding season and winter Winter territorial areas differ based on geographic location. Other populations that migrate year-round include those located in Middle America or in the western U.S. and Gulf Coast. Females typically migrate longer distances than males. These female populations located near the Great Lakes migrate nearly farther. Yearly-traveled females also migrate further than adult males, while also moving roughly the same distance as other adult females. Red-winged blackbirds migrate primarily during daytime. In general, males’ migration flocks arrive prior to females in the spring and after females in the fall.
Ecological and economic impact
According to the American ornithologist Arthur Cleveland Bent, in the northern regions of its range the eastern red-winged blackbird is almost completely beneficial from an economic perspective and there are comparatively few complaints of severe crop damage. Their diet consists almost entirely of insects, very few of which are useful species, and herb seeds. However, it causes certain damages to the grains that germinate in spring and to sweet corn in summer, while the grains are still soft, tearing the foliaceous covering of the ears and ruining them from a commercial point of view. It also attacks other grains in a limited way, but most of what it consumes is waste left in the ground. In the Midwest, where these birds are much more abundant and where cereals are grown more extensively than in the North, red-winged blackbirds and other ichterids, in late summer and fall, do great damage to grain fields, both while they are maturing as when they are harvested . However, it has been claimed that even there are beneficial because the larvae removed from corncobs and beet plants can counteract pests of caterpillars. In the southern states, they seriously harm rice by plucking seedlings in spring and eating the still-soft grains as they mature, being in this sense almost as harmful as the Bobolink. On the other hand, they are of some use in consuming weed seeds that would otherwise devalue the product.
Being one of the most numerous birds on the continent, it plays an important role in the dispersal of other species. Since red-winged blackbirds gather and rest in such large numbers, the survival of certain species that join their flocks is likely to be affected by their company. They can also be an important source of food for animals such as raccoons and mink. Likewise, populations that nest and rest in swamps could cushion the effect of predation on duck species and other animals. In summary, these birds are so numerous and active that their mere presence and natural behavior is enough to influence the environment in a visible way.
Positives: weed control and harmful insects
Through the control of insect populations through predation and unwanted herbs with the consumption of their seeds, they allow the growth of larger plants and crops. They also eat on Anthonomus grandis and Hypera postica, two species of weevils affecting cotton and alfalfa respectively, as well as harmful caterpillars of the European gypsy moth (Lymantria dispar) and of the genus Malacosoma. In some areas of the southern United States, the seeds of the common plumber (Xanthium strumarium), a weed detrimental to soybeans and cotton, seems to be an important food source for the species.
During the breeding season, the approximately 8 million red-winged blackbirds nesting in Ohio and their chicks probably consume more than of insects, an average of almost . Many of these insects, such as the weevils (Hypera spp.), Come from alfalfa fields, pastures, oat fields and other crops. In cornfields, blackbirds often feed on corn worms (Helicoverpa zea) and beetles of the genus Diabrotica. In early spring, red-winged blackbirds consume corn borers (Ostrinia nubilalis) in fields with corn stubble. However, Bendell et al. (1981) found that the economic benefit of pest control, such as larvae of that lepidopteran, by the red-winged blackbird only compensated for 20% of the damage to crops caused by this bird.
Negative aspects: consumption of cultivated grains
The red-winged blackbirds can devastate farm fields. Despite the fact that they consume weed seeds, they are known to cause great damage to agriculture due to their habits of resting in massive groups and their taste for agricultural products. may be also causes harm to plantings corn, rice, sunflower and sorghum, particularly important near roosts. Red-winged blackbird is the largest species of Icterid in North America and the most damaging to crops. From 215 birds Neotropical migrants have been identified as causing, by a wide margin, the greatest economic loss. In North America, the damage to corn crops by this species has increased since the late 1960s to early 1980s, perhaps because of the increase in the area for grain production, and due to the reduction of small areas with stubble, hayfields and uncultivated land, which, in turn, accentuated the bird's dependence on corn to ensure its livelihood.
Apparently, the male does more damage to this grain than the female. In some areas of Ohio, corn can account for up to 75% of the diet of males and only 6% of that of females in August and September. In South Dakota, in the late summer, the gizzards of the males studied contained 29% corn, while in the case of the female that number was limited to 9%.
Situation in the Midwest
The red-winged blackbird is the dominant species in the large concentrations of blackbirds that feed on the fields of sunflower, corn and small grains maturing in late summer or early fall in the Dakotas. In the 1970s, losses to sunflower and corn crops caused by blackbirds in the Dakotas exceeded $3 million annually for each case. In the northern part of the Great Plains, an area known as the Prairie Pothole Region, red-winged blackbirds are very abundant in summer. They congregate in post-reproductive flocks that significantly harm crops, particularly sunflower plantations near their home sites. Most sunflower damage occurs between mid-August and early September, when the calorie content of immature seeds is low and birds must consume more of them to satiate themselves. During this initial stage of predation on sunflower crops in which more than 75% of the total damage is caused, the red-winged blackbird represent 80% of the blackbirds observed in the fields of this seed. This period predates the massive migration of birds and most of them are of local origin. Most remain within of their native sites until the molting of their feathers is complete or nearly complete in late August or early September. Damage can be quite serious in the center and southeast of North Dakota and Northeast South Dakota, areas of high concentration of sunflower production and abundant wetlands that attract red-winged blackbird during the breeding season.
Investigations conducted between 1968 and 1979 revealed that blackbirds, notably the black-winged blackbird and the common grackle, annually destroyed less than 1% of corn crops in Ohio, amounting to a loss between 4 and 6 million dollars according to 1979 prices.
All Ohio counties experience some degree of predation on their maize crops from blackbirds, but those most affected are a few counties where marshes that roost them still abound. The counties of Ottawa, Sandusky and Lucas, on the waters of Sandusky Bay and Lake Erie, were the hardest hit. These three counties, among the 19 studied between 1968 and 1976, contained 62% of the fields in which the losses exceeded 5% and 77% of those in which they exceeded 10%. Other counties with extensive localized damage were Erie, Ashtabula also located on the coasts of the mentioned water courses and Hamilton. Almost all the plantations with damages greater than 5% were within of some important roosting of blackbirds. In the 1968–1976 period, in northeast Sandusky County and northwest Ottawa, where large roosts of up to a million birds were discovered in late summer and fall, average losses exceeded 9% in fields from the roosts, but they were less than 5% at and less than 2% at . In southwestern Ontario, in the summer of 1964, it was found that the greatest damage to the cornfields by red-winged blackbird also occurred near roosts in swamps.
Pest control
The two main options that farmers can choose from to avoid the presence of birds once corn has entered the milky stage of its maturation process are the use of the chemical 4-aminopyridine and the implementation of mechanical devices to frighten birds away. The time chosen to take measures to disperse the blackbirds is of great importance since once the birds have chosen a field to feed there they are likely to return for several days. The longer they are allowed to feed unmolested, the more difficult it will become to scare them away. Also, most of the damage is inflicted in just a few days, when the pimples are soft; consequently, control techniques will not be very useful if applied after this period.
Pest control history
As early as 1667, Massachusetts Bay settlers had enacted laws to try to reduce blackbird populations and mitigate damage to corn. According to Henry David Thoreau, a law provided that each single man in a town must kill six of those birds and, as a punishment for not doing so, he could not marry until he had complied with the aforementioned design. Obviously, since blackbirds reproduce at a much higher speed than humans marry, this control strategy was a failure. Pioneers traveling west to the Great Lakes region faced similar problems. By 1749, blackbirds were so abundant around western Lake Erie that people took turns watching over the maturing grain crops.
At the beginning of the 20th century, in some places, when the reeds dried up, these circumstances were used to kill these birds in the following way. A crew approached a roost in silence, hidden in the darkness of the night, and simultaneously lit the reeds at various points, which were quickly enveloped by a single great flame. This caused a huge tumult among the red-winged blackbird, which, lit by fire, were shot down in large numbers as they hovered in midair and screamed all over the place. Sometimes straw was used for the same purpose, which was previously scattered near reeds and alder bushes (Alnus spp.) in which they gathered to rest, the burning of which caused great consternation among the birds. The gang returned the next day to collect the hunted prey.
Arthur Cleveland Bent says that, before it was banned the sale of prey hunting in the market, red-winged blackbirds were massacred in large numbers in autumn and sold in markets. When they had put on weight on a diet of grains or rice, their small bodies were served as delicious snacks on the gourmet tables . Few could distinguish them from bobolinks (Dolichonyx oryzivorus).
In 1926, when the US Biological Survey – predecessor of the United States Fish and Wildlife Service carried out its first compilation of roosts of undesirable ichterids, it was recorded in Ohio, with its large populations of these birds and the fifth largest area allocated to the cultivation of corn among the American states, a number of complaints higher than in any other state. During the 1950s, bird control committees were organized in some counties to deal with the damage to maize caused by blackbirds and the Ohio Agricultural Experiment Station now the Ohio Agricultural Research and Development Center and the Department of Zoology and Entomology of the state university began investigating the problem.
Use of deadly traps and chemicals
Crop predation has led to the use of traps, poison and surfactants by farmers in an attempt to control populations of red-winged blackbirds; these last properties suppress waterproof feathers, making them extremely vulnerable to the cold, but their effectiveness depends on certain atmospheric conditions, namely low temperatures and rainfall. Programs in which baits were used poisoned to reduce icteride concentrations in late summer have been unsuccessful. While thousands of birds have occasionally died, the effect on large roosting-associated flocks that sometimes contain more than a million individuals is small; In addition, specimens of other species frequently die. The use of large lure traps, which often catch hundreds of birds per day, is also ineffective against these large flocks.
Use of 4-aminopyridine
4-Aminopyridine is applied to one in one hundred particles of ground corn used as bait. Generally, corn is thrown into the fields from planes that release a load of about of bait per hectare on one third of the land. Because that amount of ground corn contains around 205,000 particles, approximately 2050 toxic particles are distributed per hectare treated. The ingestion of one or more of these particles by a blackbird causes erratic flight, calls for suffering and finally death; that behavior often leads the remaining birds in the flock to leave the field. The chemical DRC-1327, which has proven useful in mitigating damage to maturing corn, operates in the same way: when a bird ingests kernels from a partially hand-peeled cob to which the chemical has been applied with a sprayer manual, his erratic flight and his pre-death calls for suffering, which span a space of between five and fifteen minutes, chase away flocks from the fields.
The initial application of 4-aminopyridine should be carried out as soon as possible after the start of the milky stage of the grain ripening process. Two other booster applications five to seven days apart are generally recommended; Often just one is sufficient, but under conditions of prolonged bird activity more than three applications may be required at shorter intervals. In Brown County (Northeast South Dakota), in 1965, hand-spread 4-aminopyridine baits at intervals of about one week reduced projected red-winged blackbird loss by maturing corn crops by 85%. The distressing behavior exhibited by the individuals affected by the chemical produced a marked fear response in other members of the flocks and the fields were free of red-winged blackbirds even when the estimated directly affected proportion was less than 1%. It has even been suggested that the continued use of 4-aminopyridine over the years could cause a change in the pattern of migration to the south, as if birds were learning to avoid areas persistently treated with the chemical.
In another experiment in the same county, the number of blackbirds making use of the treated area fell dramatically over a period of five days after treatment had begun and remained low for the remainder of the season of damage to cornfields. The results of this method were largely limited to icterides. Although common pheasant were abundant, there was no evidence that any were affected, and mortality among other bird species was negligible.
The abundance of weeds should be considered as restricting the chances that birds will find bait particles scattered on the ground, so the use of this chemical must be accompanied by a weed control program. A less obvious problem is that of the insects that remove the bait. If Gryllus are detected in a field, ground corn is expected to disappear quickly. Crickets generally select untreated particles and leave the toxic ones behind; however, the rapid decrease in the total volume of ground corn on the ground decreases the attractiveness of ground-level feeding for blackbirds. Because cricket populations are difficult to control, more frequent applications or other bird control techniques may be desirable under these circumstances. A third problem is that of heavy rains, which cover the bait with soil or drag it into cracks in the ground. Likewise, a low population density of ichterides can reduce the effectiveness of its control with 4-aminopyridine.
Non-lethal methods
A variety of devices to repel them, including electronic noise-making systems, helium balloons tied in the fields, radio-controlled aircraft, and various types of scarecrows are occasionally used in cornfields. Methods such as scarecrows, pyrotechnics, and propane cannons, may help mitigate mild predation, but only work effectively if the duration of the damage period is less than that of the birds becoming accustomed to these methods. Red-winged blackbirds quickly become accustomed to them, particularly if the crop is a prime food source in an area with few alternative sources of livelihood. Although harvesting as early as possible after the corn has dried sufficiently can limit damage by northern grazing flocks, adjusting the harvest date does not help farmers reduce losses by red-winged blackbirds during the milky stage.
One approach that has been successful in controlling roosts in the highlands is dispersal of populations through habitat alteration or bird harassment. These procedures, carried out by biologists in cooperation with local citizens, have been successful in dispersing or displacing populations of up to one million individuals. Although this dispersion can sometimes move the problem from one place to another – especially when the area has been intensively cultivated with corn and alternative food sources are not abundant, it has often been effective in solving local problem situations. Because oats and wheat grains in already harvested fields constitute an important food for blackbirds in late summer, the postponement of plowing the land with small grain stubble can lessen the pressure exerted by the predation of these birds on the maturing corn. The natural existence or planting of plants such as millet, sorghum (Sorghum spp.), Polygons (Polygonum spp.) And various grasses for example, mohas (Setaria spp.) that could be beneficial. As a general ecological principle, diversity in the types of habitats that can be maintained in regions of intense agricultural activity is related to a greater probability that the damages caused by pests are restricted to economically tolerable levels.
Studies in sweet corn fields indicate that blackbirds could often be initially attracted to maturing crops by insects. Flocks can wander the land cultivated for about a week consuming insects and weed seeds before attacking the corn. Diabrotica beetles may be especially attractive during this period. Thus the birds become habituated to feeding in the fields and quickly move from insects to corn when it enters the vulnerable milky stage. Experiments in which insect populations in sweet corn crops were treated with insecticides during the week before the grain entered the milky phase, they showed that fewer birds visited these fields and less damage was recorded to the corn in the subsequent vulnerability period than in nearby untreated land. It is probable that the abundance of weeds in cornfields also increases their attractiveness to ichterids and their control would lead to a decrease in the losses produced by these birds.
Chemical sterilization
An alternative way to reduce populations of red-winged blackbird, and therefore the damage they cause to crops, involves implementing a program that is intended to interfere with their ability to reproduce, for example, through the use of sterilizing chemicals. Due to the polygynous nature of the species, it has been suggested that such a program could become more effective if directed at males. However, the incidence of promiscuity would imply that chemical sterilization of a certain fraction of males from a local population would not result in a proportional decrease in fertile clutches.
The effects of this method would probably vary according to the type of habitat in which the treated males have established their territories. In general, there should be a higher proportion of fertile clutches in those densely populated habitats and with a greater number of renidifications (nest reconstructions). Renidifications are more common in swamps than in the highlands. In turn, individuals that reproduce in the highlands probably feed more often within their own territories than those that reproduce in swamps, which probably makes promiscuity difficult in the first type of habitat. Thus, chemical sterilization could be more effective among upland populations than among swamp populations.
Sterilizing chemicals may reduce the number of chicks produced by a successful nest. Since renidification clutches are significantly fewer in number than the original, even if a female that first mated with a sterile male then renidified and mated with a fertile male, perhaps she would produce fewer chicks than she would have had if her first clutch would have been fertile.
The number of fertile clutches in a male's territory may be limited through the sterilization of that male, but the degree to which that number decreases will depend on the disposition of the neighboring fertile territorial males and perhaps also on the number of fertile non-territorial males. in the population. Studies with artificial eggs suggest that incubation of sterile eggs would be prolonged, 98 since females normally incubate artificial eggs for around 20 days. An extended incubation of unviable eggs would result in fewer renidification attempts and, therefore, it decreased the opportunities to find a fertile mate for a female who originally mated with a sterile male.
Relationship with humans
Farmers have been known to use pesticides—such as parathion—in illegal attempts to control their populations. In the United States, such efforts are illegal because no pesticide can be used on non-target organisms, or for any use not explicitly listed on the pesticide's label. However, the USDA has deliberately poisoned this species in the past: in 2009, the Animal and Plant Health Inspection Service reported poisoning over 950,000 red-winged blackbirds in Texas and Louisiana. This poisoning has been implicated as a potential cause of the decline of the rusty blackbird, a once abundant species that has declined 99% since the 1960s and has been recently listed as Threatened on the IUCN Red List.
In the Anishinaabe languages, an indigenous language group spoken throughout much of the bird's northeastern range, this bird's names are diverse. In the Oji-Cree language, the northernmost of the Anishinaabe languages, it is called jachakanoob, while the Ojibwa language spoken in Northwestern Ontario and into Manitoba ranging immediately south of the Oji-Cree's range, the bird is called jachakanoo (with the cognates cahcahkaniw (Swampy Cree), cahcahkaluw (coastal Southern East Cree), cahcahkayuw (inland Southern East Cree), cahcahkayow (Plains Cree); the northern Algonquian languages classify the red-winged blackbird as a type of a junco or grackle, deriving the bird's name from their word for "spotted" or "marked". In the vast majority of the other Ojibwa language dialects, the bird is called memiskondinimaanganeshiinh, literally meaning "a bird with a very red damn-little shoulder-blade". However, in the Odawa language, an Anishinaabe language in southwestern Ontario and in Michigan, the bird is instead called either memeskoniinisi ("bird with a red [patch on its wing]") or memiskonigwiigaans ("[bird with a] wing of small and very red [patch]"). In N'syilxcn (Colville-Okanagan, Interior Salish language) the bird is known as ƛ̓kƛ̓aʕkək. In the Hoocąk language they are known as cooxją́ aporošucra, which describes the round red spot on its wing as well as identifying it as a blackbird.
In the Great Plains, the Lakota language, another Indigenous language spoken throughout much of the bird's range, the bird is called wabloša ("wings of red"). Its songs are described in Lakota as tōke, mat'ā nī ("oh! that I might die"), as nakun miyē ("...and me"), as miš eyā ("me too!"), and as cap'cehlī ("a beaver's running sore").
Conservation status
It is a species of least-concern. Being one of the largest and most widely distributed birds in North America, little has been done to protect it from the effects of habitat destruction and urbanization. It can survive in a wide range of environments, many populations manage to overcome the loss of natural habitats. However, red-winged blackbirds thrive in wetland areas and with the destruction of natural wetlands their population is likely to shrink. The species is protected under the Migratory Bird Treaty Act 1918, a formal treaty between the United States and Canada that was later expanded to include Mexico. This law gives them legal protection in the United States, but they can be killed "when they are found preying or about to prey on ornamental trees or trees planted for shade, crops, livestock or wildlife."
A study in Illinois indicated that red-winged blackbird populations doubled between 1908 and 1958. It had traditionally reproduced in wetlands, with Ohio primarily inhabiting swamps associated with lakes and rivers. During the 20th century, however, it adapted to man-made habitat changes and now often nests in hayfields, along roads and ditches, and elsewhere in the highlands.
Despite its successful adaptation to changes in practices related to land use, populations of red-winged blackbird have reduced the width of its range during the second half of the twentieth century, and Changes in the abundance and adequacy of grasslands have been implicated in this. In Ohio, the red-winged blackbird was negatively affected between 1966 and 1996 by the decrease in hay production, the earlier harvest of hay and the increase in crops planted in furrows, a situation homologous to that suffered by other species of birds from less numerous grasslands. Although the negative effect of the increased efficiency and the diminishing diversity of modern agricultural practices on the populations of red-winged blackbird in this state may be perceived as a positive event by the producers of corn and sunflower, the agricultural practices that have precipitated the numerical decline of this species may have caused more severe repercussions for birds from less common grasslands, such as the upland sandpiper and the Grasshopper sparrow.
In the late 1970s, Ohio hosted the highest density of blackbirds during the breeding season among all US states and Canadian provinces, However, between 1966 and 1996, the reproductive populations in this state showed a marked decrease. Decreasing the area under hay other than alfalfa between those same years was likely to reduce the availability of quality nesting habitats. Likewise, corn even though it attracts insects that red-winged blackbirds consume and despite the fact that their grains represent in themselves an energy source and soybeans which do not constitute a food source for the species, which together made up 1966 to 1996, an average of 70% of the cultivated area in Ohio, do not provide adequate nesting habitat. In turn, the large annual fluctuations in the area cultivated with hay, both from one year to the next and within the same year of a certain reproductive season, Similarly, in Ontario between 1974 and 1995, the average size of harems in a given year was positively related to hay production from the previous year, although the percentage variation in the size of harems from one season to the next the following was not linearly related to the annual percentage change in hay production; Hay production declined during this period and the size of the harems decreased from approximately three females per male to 1.6 females.
In addition in Ontario, a negative relationship was found between the North Atlantic Oscillation Index (NAO) in the six months prior to a certain reproductive season and the size of the harems. If winter mortality contributes to the decrease in the size of the harems, the annual variations in their size should be related to the annual changes in the rates of return of marked territorial males. Variations in the size of harems and male return rates per year were indeed found to be positively correlated. Although not as pronounced, the relationship between male return rates and winter NAO values was similar to that between changes in harem size and NAO.
On the other hand, in southwestern Quebec, red-winged blackbird populations doubled between 1966 and 1981, apparently in response to the development of maize production. The increased availability of residual grains in spring and summer. the reproductive season will probably play a fundamental role in the growth of this bird's population. Likewise, in the mid-1990s, the reproductive population of North Dakota which at least between 1994 and 2002 was at the forefront among all the American states and Canadian provinces in terms of population density of this species began to increase rapidly.
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197489 | https://en.wikipedia.org/wiki/SIM%20card | SIM card | A SIM (Subscriber Identity Module) card is an integrated circuit (IC) intended to securely store an international mobile subscriber identity (IMSI) number and its related key, which are used to identify and authenticate subscribers on mobile telephone devices (such as mobile phones and laptops). SIMs are also able to store address book contacts information, and may be protected using a PIN code to prevent unauthorized use.
SIMs are always used on GSM phones; for CDMA phones, they are needed only for LTE-capable handsets. SIM cards are also used in various satellite phones, smart watches, computers, or cameras. The first SIM cards were the size of credit and bank cards; sizes were reduced several times over the years, usually keeping electrical contacts the same, to fit smaller-sized devices. SIMs are transferable between different mobile devices by removing the card itself.
Technically the actual physical card is known as a universal integrated circuit card (UICC); this smart card is usually made of PVC with embedded contacts and semiconductors, with the SIM as its primary component. In practice the term "SIM card" is still used to refer to the entire unit and not simply the IC. A SIM contains a unique serial number, integrated circuit card identification (ICCID), international mobile subscriber identity (IMSI) number, security authentication and ciphering information, temporary information related to the local network, a list of the services the user has access to, and four passwords: a personal identification number (PIN) for ordinary use, and a personal unblocking key (PUK) for PIN unlocking as well as a second pair (called PIN2 and PUK2 respectively) which are used for managing fixed dialing number and some other functionality. In Europe, the serial SIM number (SSN) is also sometimes accompanied by an international article number (IAN) or a European article number (EAN) required when registering online for the subscription of a prepaid card.As of 2020, eSIM is superseding physical SIM cards in some domains, including cellular telephony. eSIM uses a software-based SIM embedded into an irremovable eUICC.
History and procurement
The SIM card is a type of smart card, the basis for which is the silicon integrated circuit (IC) chip. The idea of incorporating a silicon IC chip onto a plastic card originates from the late 1960s. Smart cards have since used MOS integrated circuit chips, along with MOS memory technologies such as flash memory and EEPROM (electrically EPROM).
The SIM was initially specified by the ETSI in the specification TS 11.11. This describes the physical and logical behaviour of the SIM. With the development of UMTS, the specification work was partially transferred to 3GPP. 3GPP is now responsible for the further development of applications like SIM (TS 51.011) and USIM (TS 31.102) and ETSI for the further development of the physical card UICC.
The first SIM card was manufactured in 1991 by Munich smart-card maker Giesecke+Devrient, who sold the first 300 SIM cards to the Finnish wireless network operator Radiolinja, who launched the world's first commercial 2G GSM cell network that year.
Today, SIM cards are considered ubiquitous, allowing over 8 billion devices to connect to cellular networks around the world daily. According to the International Card Manufacturers Association (ICMA), there were 5.4 billion SIM cards manufactured globally in 2016 creating over $6.5 billion in revenue for traditional SIM card vendors. The rise of cellular IoT and 5G networks was predicted by Ericsson to drive the growth of the addressable market for SIM cards to over 20 billion devices by 2020. The introduction of embedded-SIM (eSIM) and remote SIM provisioning (RSP) from the GSMA may disrupt the traditional SIM card ecosystem with the entrance of new players specializing in "digital" SIM card provisioning and other value-added services for mobile network operators.
Design
There are three operating voltages for SIM cards: , and (ISO/IEC 7816-3 classes A, B and C, respectively). The operating voltage of the majority of SIM cards launched before 1998 was . SIM cards produced subsequently are compatible with and . Modern cards support , and .
Modern SIM cards allow applications to load when the SIM is in use by the subscriber. These applications communicate with the handset or a server using SIM Application Toolkit, which was initially specified by 3GPP in TS 11.14. (There is an identical ETSI specification with different numbering.) ETSI and 3GPP maintain the SIM specifications. The main specifications are: ETSI TS 102 223 (the toolkit for smart cards), ETSI TS 102 241 (API), ETSI TS 102 588 (application invocation), and ETSI TS 131 111 (toolkit for more SIM-likes). SIM toolkit applications were initially written in native code using proprietary APIs. To provide interoperability of the applications, ETSI chose Java Card. A multi-company collaboration called GlobalPlatform defines some extensions on the cards, with additional APIs and features like more cryptographic security and RFID contactless use added.
Data
SIM cards store network-specific information used to authenticate and identify subscribers on the network. The most important of these are the ICCID, IMSI, authentication key (Ki), local area identity (LAI) and operator-specific emergency number. The SIM also stores other carrier-specific data such as the SMSC (Short Message service center) number, service provider name (SPN), service dialing numbers (SDN), advice-of-charge parameters and value-added service (VAS) applications. (Refer to GSM 11.11.)
SIM cards can come in various data capacities, from to at least . All can store a maximum of 250 contacts on the SIM, but while the has room for 33 Mobile country code (MCCs) or network identifiers, the version has room for 80 MNCs. This is used by network operators to store data on preferred networks, mostly used when the SIM is not in its home network but is roaming. The network operator that issued the SIM card can use this to have a phone connect to a preferred network that is more economic for the provider instead of having to pay the network operator that the phone discovered first. This does not mean that a phone containing this SIM card can connect to a maximum of only 33 or 80 networks, instead it means that the SIM card issuer can specify only up to that number of preferred networks. If a SIM is outside these preferred networks, it uses the first or best available network.
ICCID
Each SIM is internationally identified by its integrated circuit card identifier (ICCID). Nowadays ICCID numbers are also used to identify eSIM profiles, not only physical SIM cards. ICCIDs are stored in the SIM cards and are also engraved or printed on the SIM card body during a process called personalisation.
The ICCID is defined by the ITU-T recommendation E.118 as the primary account number. Its layout is based on ISO/IEC 7812. According to E.118, the number can be up to 19 digits long, including a single check digit calculated using the Luhn algorithm. However, the GSM Phase 1 defined the ICCID length as an opaque data field, 10 octets (20 digits) in length, whose structure is specific to a mobile network operator.
The number is composed of three subparts:
Issuer identification number (IIN)
Check digit
Individual account identification
Their format is as follows.
Issuer identification number (IIN)
Maximum of seven digits:
Major industry identifier (MII), 2 fixed digits, 89 for telecommunication purposes.
Country code, 2 or 3 digits, as defined by ITU-T recommendation E.164.
NANP countries, apart from Canada, use 01, i.e. prepending a zero to their common calling code +1
Canada uses 302
Russia uses 701, i.e. appending 01 to its calling code +7
Kazakhstan uses 997, even though it shares the calling code +7 with Russia
Issuer identifier, 1–4 digits.
Often identical to the Mobile country code (MCC).
Individual account identification
Its length is variable, but every number under one IIN has the same length.
Often identical to the Mobile identification number (MIN).
Check digit
Single digit calculated from the other digits using the Luhn algorithm.
With the GSM Phase 1 specification using 10 octets into which ICCID is stored as packed BCD, the data field has room for 20 digits with hexadecimal digit "F" being used as filler when necessary. In practice, this means that on GSM cards there are 20-digit (19+1) and 19-digit (18+1) ICCIDs in use, depending upon the issuer. However, a single issuer always uses the same size for its ICCIDs.
As required by E.118, the ITU-T updates a list of all current internationally assigned IIN codes in its Operational Bulletins which are published twice a month (the last as of January 2019 was No. 1163 from 1 January 2019). ITU-T also publishes complete lists: as of August 2023, the list issued on 1 December 2018 was current, having all issuer identifier numbers before 1 December 2018.
International mobile subscriber identity (IMSI)
SIM cards are identified on their individual operator networks by a unique international mobile subscriber identity (IMSI). Mobile network operators connect mobile phone calls and communicate with their market SIM cards using their IMSIs. The format is:
The first three digits represent the Mobile country code (MCC).
The next two or three digits represent the Mobile network code (MNC). Three-digit MNC codes are allowed by E.212 but are mainly used in the United States and Canada. One MCC can have both 2 digit and 3 digit MNCs, an example is 350 007.
The next digits represent the Mobile identification number (MSIN).
Normally there are 10 digits, but can be fewer in the case of a 3-digit MNC or if national regulations indicate that the total length of the IMSI should be less than 15 digits.
Digits are different from country to country.
Authentication key (Ki)
The Ki is a 128-bit value used in authenticating the SIMs on a GSM mobile network (for USIM network, the K is still needed but other parameters are also needed). Each SIM holds a unique Ki assigned to it by the operator during the personalisation process. The Ki is also stored in a database (termed authentication center or AuC) on the carrier's network.
The SIM card is designed to prevent someone from getting the Ki by using the smart-card interface. Instead, the SIM card provides a function, Run GSM Algorithm, that the phone uses to pass data to the SIM card to be signed with the Ki. This, by design, makes using the SIM card mandatory unless the Ki can be extracted from the SIM card, or the carrier is willing to reveal the Ki. In practice, the GSM cryptographic algorithm for computing a signed response (SRES_1/SRES_2: see steps 3 and 4, below) from the Ki has certain vulnerabilities that can allow the extraction of the Ki from a SIM card and the making of a duplicate SIM card.
Authentication process:
When the mobile equipment starts up, it obtains the international mobile subscriber identity (IMSI) from the SIM card, and passes this to the mobile operator, requesting access and authentication. The mobile equipment may have to pass a PIN to the SIM card before the SIM card reveals this information.
The operator network searches its database for the incoming IMSI and its associated Ki.
The operator network then generates a random number (RAND, which is a nonce) and signs it with the Ki associated with the IMSI (and stored on the SIM card), computing another number, that is split into the Signed Response 1 (SRES_1, 32 bits) and the encryption key Kc (64 bits).
The operator network then sends the RAND to the mobile equipment, which passes it to the SIM card. The SIM card signs it with its Ki, producing Signed Response 2 (SRES_2) and Kc, which it gives to the mobile equipment. The mobile equipment passes SRES_2 on to the operator network.
The operator network then compares its computed SRES_1 with the computed SRES_2 that the mobile equipment returned. If the two numbers match, the SIM is authenticated and the mobile equipment is granted access to the operator's network. Kc is used to encrypt all further communications between the mobile equipment and the operator.
Location area identity
The SIM stores network state information, which is received from the location area identity (LAI). Operator networks are divided into location areas, each having a unique LAI number. When the device changes locations, it stores the new LAI to the SIM and sends it back to the operator network with its new location. If the device is power cycled, it takes data off the SIM, and searches for the prior LAI.
SMS messages and contacts
Most SIM cards store a number of SMS messages and phone book contacts. It stores the contacts in simple "name and number" pairs. Entries that contain multiple phone numbers and additional phone numbers are usually not stored on the SIM card. When a user tries to copy such entries to a SIM, the handset's software breaks them into multiple entries, discarding information that is not a phone number. The number of contacts and messages stored depends on the SIM; early models stored as few as five messages and 20 contacts, while modern SIM cards can usually store over 250 contacts.
Formats
SIM cards have been made smaller over the years; functionality is independent of format. Full-size SIM was followed by mini-SIM, micro-SIM, and nano-SIM. SIM cards are also made to embed in devices.
All versions of the non-embedded SIM cards share the same ISO/IEC 7816 pin arrangement.
Full-size SIM
The full-size SIM (or 1FF, 1st form factor) was the first form factor to appear. It was the size of a credit card (85.60 mm × 53.98 mm × 0.76 mm).
Mini-SIM
The mini-SIM (or 2FF) card has the same contact arrangement as the full-size SIM card and is normally supplied within a full-size card carrier, attached by a number of linking pieces. This arrangement (defined in ISO/IEC 7810 as ID-1/000) lets such a card be used in a device that requires a full-size card or in a device that requires a mini-SIM card, after breaking the linking pieces. As the full-size SIM is obsolete, some suppliers refer to the mini-SIM as a "standard SIM" or "regular SIM".
Micro-SIM
The micro-SIM (or 3FF) card has the same thickness and contact arrangements, but reduced length and width as shown in the table above.
The micro-SIM was introduced by the European Telecommunications Standards Institute (ETSI) along with SCP, 3GPP (UTRAN/GERAN), 3GPP2 (CDMA2000), ARIB, GSM Association (GSMA SCaG and GSMNA), GlobalPlatform, Liberty Alliance, and the Open Mobile Alliance (OMA) for the purpose of fitting into devices too small for a mini-SIM card.
The form factor was mentioned in the December 1998 3GPP SMG9 UMTS Working Party, which is the standards-setting body for GSM SIM cards, and the form factor was agreed upon in late 2003.
The micro-SIM was designed for backward compatibility. The major issue for backward compatibility was the contact area of the chip. Retaining the same contact area makes the micro-SIM compatible with the prior, larger SIM readers through the use of plastic cutout surrounds. The SIM was also designed to run at the same speed (5 MHz) as the prior version. The same size and positions of pins resulted in numerous "How-to" tutorials and YouTube videos with detailed instructions how to cut a mini-SIM card to micro-SIM size.
The chairman of EP SCP, Klaus Vedder, said
Micro-SIM cards were introduced by various mobile service providers for the launch of the original iPad, and later for smartphones, from April 2010. The iPhone 4 was the first smartphone to use a micro-SIM card in June 2010, followed by many others.
Nano-SIM
After a debate in early 2012 between a few designs created by Apple, Nokia and RIM, Apple's design for an even smaller SIM card was accepted by the ETSI. The nano-SIM (or 4FF) card was introduced in June 2012, when mobile service providers in various countries first supplied it for phones that supported the format. The nano-SIM measures and reduces the previous format to the contact area while maintaining the existing contact arrangements. A small rim of isolating material is left around the contact area to avoid short circuits with the socket. The nano-SIM can be put into adapters for use with devices designed for 2FF or 3FF SIMs, and is made thinner for that purpose, and telephone companies give due warning about this. 4FF is thick, compared to the of its predecessors.
The iPhone 5, released in September 2012, was the first device to use a nano-SIM card, followed by other handsets.
Security
In July 2013, Karsten Nohl, a security researcher from SRLabs, described vulnerabilities in some SIM cards that supported DES, which, despite its age, is still used by some operators. The attack could lead to the phone being remotely cloned or let someone steal payment credentials from the SIM. Further details of the research were provided at BlackHat on 31 July 2013. In response, the International Telecommunication Union said that the development was "hugely significant" and that it would be contacting its members.
In February 2015, The Intercept reported that the NSA and GCHQ had stolen the encryption keys (Ki's) used by Gemalto (now known as Thales DIS, manufacturer of 2 billion SIM cards annually) ), enabling these intelligence agencies to monitor voice and data communications without the knowledge or approval of cellular network providers or judicial oversight. Having finished its investigation, Gemalto claimed that it has “reasonable grounds” to believe that the NSA and GCHQ carried out an operation to hack its network in 2010 and 2011, but says the number of possibly stolen keys would not have been massive.
In September 2019, Cathal Mc Daid, a security researcher from Adaptive Mobile Security, described how vulnerabilities in some SIM cards that contained the S@T Browser library were being actively exploited. This vulnerability was named Simjacker. Attackers were using the vulnerability to track the location of thousands of mobile phone users in several countries. Further details of the research were provided at VirusBulletin on 3 October 2019.
Developments
When GSM was already in use, the specifications were further developed and enhanced with functionality such as SMS and GPRS. These development steps are referred as releases by ETSI. Within these development cycles, the SIM specification was enhanced as well: new voltage classes, formats and files were introduced.
USIM
In GSM-only times, the SIM consisted of the hardware and the software. With the advent of UMTS, this naming was split: the SIM was now an application and hence only software. The hardware part was called UICC. This split was necessary because UMTS introduced a new application, the universal subscriber identity module (USIM). The USIM brought, among other things, security improvements like mutual authentication and longer encryption keys, and an improved address book.
UICC
"SIM cards" in developed countries today are usually UICCs containing at least a SIM application and a USIM application. This configuration is necessary because older GSM only handsets are solely compatible with the SIM application and some UMTS security enhancements rely on the USIM application.
Other variants
On cdmaOne networks, the equivalent of the SIM card is the R-UIM and the equivalent of the SIM application is the CSIM.
A virtual SIM is a mobile phone number provided by a mobile network operator that does not require a SIM card to connect phone calls to a user's mobile phone.
Embedded SIM (eSIM)
An embedded SIM (eSIM) is a form of programmable SIM that is embedded directly into a device. The surface mount format provides the same electrical interface as the full size, 2FF and 3FF SIM cards, but is soldered to a circuit board as part of the manufacturing process. In M2M applications where there is no requirement to change the SIM card, this avoids the requirement for a connector, improving reliability and security. An eSIM can be provisioned remotely; end-users can add or remove operators without the need to physically swap a SIM from the device or use multiple eSIM profiles at the same time.
The eSIM standard, initially introduced in 2016, has progressively supplanted traditional physical SIM cards across various sectors, notably in cellular telephony. In September 2017, Apple introduced the Apple Watch Series 3 featuring eSIM. In October 2018, Apple introduced the iPad Pro (3rd generation), which was the first iPad to support eSIM. In September 2022, Apple introduced the iPhone 14 series which was the first eSIM exclusive iPhone in the United States.
Integrated SIM (iSIM)
An integrated SIM (iSIM) is a form of SIM directly integrated into the modem chip or main processor of the device itself. As a consequence they are smaller, cheaper and more reliable than eSIMs, they can improve security and ease the logistics and production of small devices i.e. for IoT applications. In 2021, Deutsche Telekom introduced the nuSIM, an "Integrated SIM for IoT".
Usage in mobile phone standards
The use of SIM cards is mandatory in GSM devices.
The satellite phone networks Iridium, Thuraya and Inmarsat's BGAN also use SIM cards. Sometimes, these SIM cards work in regular GSM phones and also allow GSM customers to roam in satellite networks by using their own SIM cards in a satellite phone.
Japan's 2G PDC system (which was shut down in 2012; SoftBank Mobile shut down PDC from 31 March 2010) also specified a SIM, but this has never been implemented commercially. The specification of the interface between the Mobile Equipment and the SIM is given in the RCR STD-27 annexe 4. The Subscriber Identity Module Expert Group was a committee of specialists assembled by the European Telecommunications Standards Institute (ETSI) to draw up the specifications (GSM 11.11) for interfacing between smart cards and mobile telephones. In 1994, the name SIMEG was changed to SMG9.
Japan's current and next-generation cellular systems are based on W-CDMA (UMTS) and CDMA2000 and all use SIM cards. However, Japanese CDMA2000-based phones are locked to the R-UIM they are associated with and thus, the cards are not interchangeable with other Japanese CDMA2000 handsets (though they may be inserted into GSM/WCDMA handsets for roaming purposes outside Japan).
CDMA-based devices originally did not use a removable card, and the service for these phones is bound to a unique identifier contained in the handset itself. This is most prevalent in operators in the Americas. The first publication of the TIA-820 standard (also known as 3GPP2 C.S0023) in 2000 defined the Removable User Identity Module (R-UIM). Card-based CDMA devices are most prevalent in Asia.
The equivalent of a SIM in UMTS is called the universal integrated circuit card (UICC), which runs a USIM application. The UICC is still colloquially called a SIM card.
SIM and carriers
The SIM card introduced a new and significant business opportunity for who lease capacity from one of the network operators rather than owning or operating a cellular telecoms network and only provide a SIM card to their customers. MVNOs first appeared in Denmark, Hong Kong, Finland and the UK. By 2011 they existed in over 50 countries, including most of Europe, the United States, Canada, Mexico, Australia and parts of Asia, and accounted for approximately 10% of all mobile phone subscribers around the world.
On some networks, the mobile phone is locked to its carrier SIM card, meaning that the phone only works with SIM cards from the specific carrier. This is more common in markets where mobile phones are heavily subsidised by the carriers, and the business model depends on the customer staying with the service provider for a minimum term (typically 12, 18 or 24 months). SIM cards that are issued by providers with an associated contract, but where the carrier does not provide a mobile device (such as a mobile phone) are called SIM-only deals. Common examples are the GSM networks in the United States, Canada, Australia, and Poland. UK mobile networks ended SIM lock practices in December 2021. Many businesses offer the ability to remove the SIM lock from a phone, effectively making it possible to then use the phone on any network by inserting a different SIM card. Mostly, GSM and 3G mobile handsets can easily be unlocked and used on any suitable network with any SIM card.
In countries where the phones are not subsidised, e.g., India, Israel and Belgium, all phones are unlocked. Where the phone is not locked to its SIM card, the users can easily switch networks by simply replacing the SIM card of one network with that of another while using only one phone. This is typical, for example, among users who may want to optimise their carrier's traffic by different tariffs to different friends on different networks, or when travelling internationally.
In 2016, carriers started using the concept of automatic SIM reactivation whereby they let users reuse expired SIM cards instead of purchasing new ones when they wish to re-subscribe to that operator. This is particularly useful in countries where prepaid calls dominate and where competition drives high churn rates, as users had to return to a carrier shop to purchase a new SIM each time they wanted to churn back to an operator.
SIM-only
Commonly sold as a product by mobile telecommunications companies, "SIM-only" refers to a type of legally liability contract between a mobile network provider and a customer. The contract itself takes the form of a credit agreement and is subject to a credit check.
SIM-only contracts can be pre-pay - where the subscriber buys credit before use (often called pay as you go, abbreviated to PAYG), or post-pay, where the subscriber pays in arrears, typically monthly.
Within a SIM-only contract, the mobile network provider supplies their customer with just one piece of hardware, a SIM card, which includes an agreed amount of network usage in exchange for a monthly payment. Network usage within a SIM-only contract can be measured in minutes, text, data or any combination of these. The duration of a SIM-only contract varies depending on the deal selected by the customer, but in the UK they are typically available over 1, 3, 6, 12 or 24-month periods.
SIM-only contracts differ from mobile phone contracts in that they do not include any hardware other than a SIM card. In terms of network usage, SIM-only is typically more cost-effective than other contracts because the provider does not charge more to offset the cost of a mobile device over the contract period. The short contract length is one of the key features of SIM-only made possible by the absence of a mobile device.
SIM-only is increasing in popularity very quickly. In 2010 pay monthly based mobile phone subscriptions grew from 41 percent to 49 percent of all UK mobile phone subscriptions. According to German research company GfK, 250,000 SIM-only mobile contracts were taken up in the UK during July 2012 alone, the highest figure since GfK began keeping records.
Increasing smartphone penetration combined with financial concerns is leading customers to save money by moving onto a SIM-only when their initial contract term is over.
Multiple-SIM devices
Dual SIM devices have two SIM card slots for the use of two SIM cards, from one or multiple carriers. Multiple SIM devices are commonplace in developing markets such as in Africa, East Asia, South Asia and Southeast Asia, where variable billing rates, network coverage and speed make it desirable for consumers to use multiple SIMs from competing networks. Dual-SIM phones are also useful to separate one's personal phone number from a business phone number, without having to carry multiple devices. Some popular devices, such as the BlackBerry KeyOne, have dual-SIM variants; however, dual-SIM devices were not common in the US or Europe due to lack of demand. This has changed with mainline products from Apple and Google featuring either two SIM slots or a combination of a physical SIM slot and an eSIM.
In September 2018, Apple introduced iPhone XS, iPhone XS Max, and iPhone XR featuring Dual SIM (nano-SIM and eSIM) and Apple Watch Series 4 featuring Dual eSIM.
Thin SIM
A thin SIM (or overlay SIM or SIM overlay) is a very thin device shaped like a SIM card, approximately 120 microns ( inch) thick. It has contacts on its front and back. It is used by placing it on top of a regular SIM card. It provides its own functionality while passing through the functionality of the SIM card underneath. It can be used to bypass the mobile operating network and run custom applications, particularly on non-programmable cell phones.
Its top surface is a connector that connects to the phone in place of the normal SIM. Its bottom surface is a connector that connects to the SIM in place of the phone. With electronics, it can modify signals in either direction, thus presenting a modified SIM to the phone, and/or presenting a modified phone to the SIM. (It is a similar concept to the Game Genie, which connects between a game console and a game cartridge, creating a modified game). Similar devices have also been developed for iPhones to circumvent SIM card restrictions on carrier-locked models.
In 2014, Equitel, an MVNO operated by Kenya's Equity Bank, announced its intention to begin issuing thin SIMs to customers, raising security concerns by competition, particularly concerning the safety of mobile money accounts. However, after months of security testing and legal hearings before the country's Parliamentary Committee on Energy, Information and Communications, the Communications Authority of Kenya (CAK) gave the bank the green light to roll out its thin SIM cards.
| Technology | Telecommunications | null |
197590 | https://en.wikipedia.org/wiki/Madrid%20Metro | Madrid Metro | The Madrid Metro (Spanish: Metro de Madrid) is a rapid transit system serving the city of Madrid, capital of Spain. The system is the 14th longest rapid transit system in the world, with a total length of 293 km (182 mi). Its growth between 1995 and 2007 put it among the fastest-growing networks in the world at the time. However, the European debt crisis greatly slowed expansion plans, with many projects being postponed and canceled. Unlike normal Spanish road and rail traffic, which drive on the right, Madrid Metro trains use left-hand running on all lines because traffic in Madrid drove on the left until 1924, five years after the system started operating.
Trains are in circulation every day from 6:00 am until 1:30 am, though during the weekends, this schedule was to be extended by one more hour in the morning in 2020. Furthermore, the regional government intended to keep stations opened around the clock during these days from 2023 onwards. It had only stayed open 24 hours during the 2017 World Pride and during the 2021 Madrid snowstorm.
A light rail system feeding the metro opened in 2007 called Metro Ligero ("light metro"). The Cercanías system works in conjunction with the metro, with a majority of its stations providing access to the underground network.
As of January 2024, the Madrid Metro has 1,710 escalators and 559 elevators.
History
1916–1918: conception and financing
On 19 September 1916, a royal decree approved the 4-line plan for the creation of the metro of Madrid. The engineers who created the plan - Mendoza, González Echarte, and Otamendi - then began the process of raising 8 million pesetas to begin the first phase of the project, the construction of Line 1 from Sol to Cuatro Caminos. Carlos Mendoza made contact with Enrique Ocharán, the director of Banco de Vizcaya, who offered 4 million pesetas on the condition that the public pledged an additional 4 million.
Mengemor published a brochure to persuade people to make donations. The men were able to raise 2.5 million pesetas of the 4 million they needed. King Alfonso XIII intervened and invested 1.45 million pesetas of his own money.
1919: construction and inauguration
The first phase of construction was finished in 1919. It was constructed in a narrow section and the stations had platforms. The enlargement of this line and the construction of two others followed shortly after 1919. The Madrid metro was inaugurated on 17 October 1919 by King Alfonso XIII. At the time of inauguration, the metro had just one line, which ran for between Puerta del Sol and Cuatro Caminos, with eight stops.
The king, the royal family, and others took part in the first official metro ride which went from Cuatro Caminos to Ríos Rosas and took 40 seconds. There they stopped for one minute, before traveling to the Chamberí station which took 45 seconds. The trip went all the way to the end point, Sol. The king and his family then rode the metro back to Cuatro Caminos from Sol, this time without stopping. The journey took 7 minutes and 46 seconds.
After the journey, a lunch was served on the Cuatro Caminos platform, and the engineers were congratulated for creating a "miracle."
Two days later, on 19 October 1919, the Madrid metro was opened to the public. On its first day, 390 trains ran, 56,220 passengers rode the metro, and the company earned 8,433 pesetas from ticket fares.
During November and December 1919, the metro had an average of 43,537 passengers a day and earned an average of 6,530 pesetas a day from ticket sales. Due to their success, the company decided to expand more, and created 12,000 new shares to sell to the public to raise more funds to fund further expansion.
1920–1921: expansion of Line 1 and construction of Line 2
The Company then began to gather materials necessary to expand the Line 1 from Sol with the new stations Progreso, Antón Martín and finally Atocha. The latter was then and is now an important train station for mainline rail.
On 31 July 1920 the company submitted it proposal to extend Line 1 from Atocha to Puente de Vallecas. In 1921 the company declared its interest in beginning the line from Sol to Ventas, with the first phase of the project being built from Sol to Goya, along Calle Alcalá.
Work began on 27 March 1921 to expand the Line 1 from Atocha to Vallecas, and to begin construction on a line from Sol to Goya.
On 26 December 1921 the Sol-Atocha section of the Line 1 was inaugurated, adding three new metro stops to the line: Progreso, Antón Martín, and Atocha. The king and queen, Don Alfonso XIII and Doña Victoria, attended the inauguration.
1922 and onwards
In 1924, traffic in Madrid switched from driving on the left to driving on the right, but the lines of the Madrid Metro kept operating on the left hand side. In 1936, the network had three lines and a branch line between Ópera and the old Estación del Norte (now Príncipe Pío). All these stations served as air raid shelters during the Spanish Civil War. After the Civil war, the public works to extend the network went on little by little. In 1944, a fourth line was constructed, absorbing the branch of Line 2 between Goya and Diego de León in 1958, a branch that had been intended to be part of Line 4 since its construction but was operated as a branch of Line 2 until construction works had finished.
In the 1960s, a suburban railway was constructed between Plaza de España and Carabanchel, linked to lines 2 (at Noviciado station with a long transfer) and 3. A fifth metro line was constructed as well with narrow sections, but 90 m platforms. Shortly after opening the first section of Line 5, the platforms of Line 1 were enlarged from 60 to 90 m, permanently closing Chamberí station since it was too close to Iglesia (less than 500 m). Chamberí has been closed ever since and was recently reopened as a museum.
In the early 1970s, the network was greatly expanded to cope with the influx of population and urban sprawl from Madrid's economic boom. New lines were planned with larger 115 m long platforms. Lines 4 and 5 were enlarged as well. In 1979, bad management led to a crisis. Projects that had already started were finished during the 1980s and all remaining ones were abandoned. After all those projects, of rail track was completed by 1983 and the suburban railway had also disappeared since it had been extended to Alonso Martínez and subsequently converted to the new Line 10.
Expansion from the 1990s
Work on a major expansion of the metro began in 1995, with of new line and 132 new stations opened by 2011, built in 4 phases. The average construction pace throughout that era (more than of new line per year) was among the fastest in the world at that time and was equalled or surpassed by only very few metros in the global north (one notable example being Seoul Metro) either then or since.
This included the extension of lines 1, 4 and 7 and the construction of a new Line 11 towards the outlying areas of Madrid. Lines 8 and 10 were joined into a longer Line 10 and a new Line 8 was constructed to expand the underground network towards the airport. The enlarged Line 9 was the first to leave the outskirts of Madrid to arrive in Rivas-Vaciamadrid and Arganda del Rey, two satellite towns located in the southeast of Madrid. Control of the network was transferred to a public enterprise, Metro de Madrid S.A.
In the early 2000s, a huge project installed approximately of new metro tunnels. This construction included a direct connection between downtown Madrid (Nuevos Ministerios) and the airport, a further extension of Line 8, and adding service to the outskirts with a 40 km loop called MetroSur serving Madrid's southern suburbs.
MetroSur, one of the largest ever civil engineering projects in Europe, opened on 11 April 2003. It included of tunnel and 28 new stations, with a new interchange station on Line 10, connecting it to the city centre and stations linking to the local train network. Its construction began in June 2000 and the whole loop was completed in less than three years. It connects Getafe, Móstoles, Alcorcón, Fuenlabrada, and Leganés, five towns located in the south of Madrid. As the metro line is part of a project to develop the area, some stations lay in sparsely populated places or were even surrounded by fields at the time of opening.
Most of the efforts of Madrid regional government in 2000s were channeled towards the enlargement of the Metro network. In the 2003–2007 term, President Esperanza Aguirre funded a multibillion-euro project, which added new lines, and joined or extended almost all of the existing metro lines. The project included the addition of of railway and the construction of 80 new stations. It brought stations to many districts that had never previously had Metro service (Villaverde, Manoteras, Carabanchel Alto, La Elipa, Pinar de Chamartín) and to the eastern and northern outskirts as well (Coslada, San Fernando de Henares, Alcobendas, San Sebastián de los Reyes). For the first time in Madrid, three interurban light rails (Metro Ligero or ML) lines were built to the western outskirts (Pozuelo de Alarcón, Boadilla del Monte) – mL2 and mL3 – and to the new northern districts of Sanchinarro and Las Tablas – mL1. As a last minute addition, a project on line 8 connected it to the new T4 terminal of Madrid-Barajas Airport.
The Great Recession and subsequent European debt crisis, which affected Spain particularly hard, brought metro expansion to a halt as austerity measures all but dried up available funding.
Station design and setup
The age of many Madrid Metro stations is evident by their design: Older stations on the narrow lines are often quite compact, similar to stations of the Paris Metro. They were decorated with tilings in different colour schemes depending on the station similar to the scheme architect Alfred Grenander implemented at Berlin U-Bahn around the same time. In recent years, most of these stations have been refurbished with single-coloured plates matching those in the newest ones, alongside other improvements such as modernized signalling technology and elevators. The stations built between the late 70s and the early 90s are slightly more spacious, with most of them having cream-coloured walls.
On the other hand, the most recent stations are built with space in mind, and have natural-like lighting and ample entryways. The colour scheme varies between stations, using single-colored plates and covering the whole station in light colors. Recently built transfer stations have white walls, but this is not the norm.
Most stations are built with two side platforms, but a handful of them (the busiest transfers) have a central island platform in addition to the side platforms theoretically dedicated to exits. This system was originally used on the Barcelona Metro and is called the Spanish solution. The 12 stations with this setup are:
Miguel Hernández
Campamento, Carabanchel
Avenida de América, Manuel Becerra, Sainz de Baranda, Pacífico, Plaza Elíptica, Oporto, Laguna
Avenida de América, Pueblo Nuevo
Two stations have cross-platform interchange arrangement with two island platforms, which allows extremely fast transfers between lines. Both of these stations are on Line , with cross-platform interchanges at Príncipe Pío (with ) and Casa de Campo (with ). On both occasions, Line 10 uses the outside tracks, so passengers unboarding there leave through the "right" side of the train instead of the usual left side.
In addition, 10 stations are built with just one island platform instead of the usual side platforms. These stations are:
Pinar de Chamartín (one island platform and two side platforms)
Almendrales, Villaverde Alto
Aluche
Feria de Madrid, Aeropuerto T4
Rivas Urbanizaciones, Arganda del Rey
Joaquín Vilumbrales
Another system is where there is one island platform with one side platform. This system is used in two stations on lines 2 and 4 as termini, and three stations on Lines 7, 9, and 10 where it is required for passengers to change to smaller trains to continue their journeys, normally to towns outside Madrid like Alcobendas or Coslada. This is done so the island platform can be used for passengers to change easily between trains. In the latter three stations, the island platform is equipped with fare gates for tickets to be validated for travel between fare zones. These stations are:
Cuatro Caminos
Argüelles
Estadio Metropolitano
Puerta de Arganda
Tres Olivos
Overhead power supply system
Since 1999 Metro de Madrid has used a patented system for its installations: a solid rail hung from the ceiling of the tunnels, instead of the usual copper or aluminium wire hung from overhead gantries at regular intervals. This type of overhead line is rigid, making it more robust and less prone to failures. Installations outside tunnels are rare, as they require many more support structures compared to traditional wire based overhead lines, making them more expensive to install. This system of rigid overhead power supply is also used in other metro systems. Similar installations exist in some mainline rail tunnels where space is limited, e.g. Leipzig City Tunnel or the lower level of Berlin Hauptbahnhof,
Lines
The Metro conventional network has 242 stations on 12 lines plus one branch line, totalling , of which approximately 96% of stations are underground. The only surface parts are between Empalme and west of Eugenia de Montijo (); between Lago and north of Casa de Campo (); and between south of Puerta de Arganda and Arganda del Rey (), for a total of 8 aboveground stations. Additionally, some of Metro Ligero (light rail) lines across serve the various regions of the metropolitan area which have been deemed not populated enough to justify the extraordinary spending of new Metro lines. Combined, they have 38 stops, of which 4 also connect to the conventional Metro system. Most of the ML track length is on surface, usually running on platforms separated from normal road traffic. However, ML1 line has some underground stretches and stations.
Traditionally, the Madrid metro was restricted to the city proper, but today nearly one third of its track length runs outside the border of the Madrid municipality. Today, the Metro network is divided in six regions:
MetroMadrid (zone A): the core network inside the Madrid city borders, with over two-thirds of the overall length. Also includes light rail line 1.
MetroSur (zones B1 and B2): line 12 and the last two stations of line 10, Joaquín Vilumbrales and Puerta del Sur. Runs through the southern cities of Alcorcón, Leganés, Getafe, Fuenlabrada and Móstoles.
MetroEste (zone B1): a prolongation of line 7 from Estadio Metropolitano to Hospital del Henares through the municipalities of Coslada and San Fernando de Henares.
MetroNorte (zone B1): opened in 2007, includes the stretch of line 10 from La Granja to Hospital Infanta Sofía. Serves the northern outskirts of Madrid and the towns of Alcobendas and San Sebastián de los Reyes. There is a train interchange inside the line at Tres Olivos station.
MLO (zones B1 and B2): comprised by light rail lines 2 and 3. Connects the towns of Pozuelo de Alarcón and Boadilla del Monte to line 10 at Colonia Jardín station.
TFM (zones B1, B2 and B3): a prolongation of line 9 from Puerta de Arganda, the first ever outside the borders of Madrid, services the cities of Rivas-Vacíamadrid and Arganda del Rey.
At most of the borders between the regions, one has to switch trains even when staying in the same line, because the train frequency is higher in the core MetroMadrid than in the outer regions.
Madrid also has an extensive commuter train (Cercanías) network operated by Renfe, the national rail line, which is intermodal with the metro network. In fact, 22 Cercanías stations have connections to the Metro network, which is indicated on the official map by the Cercanías logo. Many of the new lines since 1999 have been built to link to or end at Cercanías stations, like the ML2 line, which ends at the Aravaca station providing a fast entry into Madrid though the C-7 or C-10 commuter lines and arriving in only one step to the bus and Metro hub Príncipe Pío ( ).
| Technology | Spain | null |
197742 | https://en.wikipedia.org/wiki/Sambucus | Sambucus | Sambucus is a genus of flowering plants in the family Adoxaceae. The various species are commonly referred to as elder, elderflower or elderberry.
Description
The oppositely arranged leaves are pinnate with 5–9 leaflets (or, rarely, 3 or 11). Each leaf is long, and the leaflets have serrated margins. They bear large clusters of small white or cream-colored flowers in late spring; these are followed by clusters of small black, blue-black, or red berries (rarely yellow or white).
Chemistry
Sambucus fruit is rich in anthocyanidins that combine to give elderberry juice an intense blue-purple coloration that turns reddish on dilution with water.
Taxonomy
The taxonomy of the genus Sambucus L., originally described by Carl Linnaeus and hence its botanical authority, has been complicated by its wide geographical distribution and morphological diversity. This has led to overdescription of the species and infraspecific taxa (subspecies, varieties or forms).
The genus was formerly placed in the honeysuckle family, Caprifoliaceae, but was reclassified as Adoxaceae due to genetic and morphological comparisons to plants in the genus Adoxa.
Species recognized in this genus are:
– Himalaya and eastern Asia
Sambucus australasica – New Guinea, eastern Australia
Sambucus australis – South America
Sambucus canadensis – eastern North America
Sambucus cerulea – western North America
Sambucus ebulus – central and southern Europe, northwest Africa and southwest Asia
Sambucus gaudichaudiana – south eastern Australia
Sambucus javanica – southeastern Asia
Sambucus lanceolata – Madeira Island
– Korea, southeast Siberia
– western North America
– southwest North America
Sambucus nigra – Europe and North America
Sambucus orbiculata – western North America
Sambucus palmensis – Canary Islands
Sambucus peruviana – Costa Rica, Panama and northwest South America
Sambucus pubens – northern North America
Sambucus racemosa – northern, central and southeastern Europe, northwest Asia, western North America
– eastern Asia
Sambucus sieboldiana – Japan and Korea
– southeastern United States
Sambucus tigranii – southwest Asia
Sambucus velutina – southwestern North America
– western Himalayas
– northeast Asia
Etymology
The name comes from the Ancient Greek word (), an ancient wind instrument, about the removal of pith from the twigs to make whistles.
Distribution and habitat
The genus occurs in temperate to subtropical regions of the world. More widespread in the Northern Hemisphere, its Southern Hemisphere occurrence is restricted to parts of Australasia and South America. Many species are widely cultivated for their ornamental leaves, flowers, and fruit.
Elder commonly grows near farms and homesteads. It is a nitrogen-dependent plant and thus is generally found near places of organic waste disposal. Elders are often grown as a hedgerow plant in Britain since they take very fast, can be bent into shape easily, and grow quite profusely, thus having gained the reputation of being 'an instant hedge'. It is not generally affected by soil type or pH level and will virtually grow anywhere sufficient sunlight is available.
Ecology
S. callicarpa berries are consumed by birds and mammals. In Northern California, elderberries are eaten by migrating band-tailed pigeons. Elders are used as food plants by the larvae of some Lepidoptera species including brown-tail, buff ermine, dot moth, emperor moth, engrailed moth, swallow-tailed moth and the . The crushed foliage and immature fruit have a strong fetid smell. Valley elderberry longhorn beetles in California are very often found around red or blue elderberry bushes. Females lay their eggs on the bark. Strong-scented flowers in wild populations of European elder (S. nigra) attract numerous, minute flower thrips which may contribute to the transfer of pollen between inflorescences.
Cultivation
Traditional uses of Sambucus involved berries, seeds, leaves, and flowers or component extracts. Ornamental varieties of Sambucus are grown in gardens for their showy flowers, fruits and lacy foliage which support habitat for wildlife. Of the many native species, three are used as ornamentals: S. canadensis, S. nigra, and S. racemosa.
Toxicity
The uncooked berries and other parts of plants from this genus are poisonous. Leaves, twigs, branches, seeds, roots, flowers, and berries of Sambucus plants produce cyanogenic glycosides, which have toxic properties. Ingesting a sufficient quantity of cyanogenic glycosides from berry juice, flower tea, or beverages made from fresh leaves, branches, and fruit has been shown to cause illness, including nausea, vomiting, abdominal cramps, diarrhea, and weakness. In August 1983, a group of 25 people in Monterey County, California, became ill after ingesting elderberry juice pressed from fresh, uncooked S. mexicana berries, leaves, and stems. The concentration of cyanogenic glycosides is higher in tea made from flowers (or leaves) than from the berries.
The seeds of S. callicarpa are reported to be poisonous and may cause vomiting or diarrhea.
Uses
The cooked berries (pulp and skin) of most species of Sambucus are edible.
Nutrition
Raw elderberries are 80% water, 18% carbohydrates, and less than 1% each of protein and fat. In a amount, elderberries supply of food energy and are a rich source of vitamin C, providing 43% of the Daily Value (DV). Elderberries also have moderate contents of vitamin B6 (18% DV) and iron (12% DV), with no other nutrients in significant content.
Dietary supplement
Elderberry fruit or flowers are used as dietary supplements to prevent or provide relief from minor diseases, such as flu, colds, constipation, and other conditions, served as a tea, extract or in a capsule. The use of elderberry supplements increased early in the COVID-19 pandemic. There is insufficient research to establish its effectiveness for such uses, or its safety profile. The raw or unripe fruit of S. nigra or its extracts may contain a cyanogenic glycoside that is potentially toxic.
Traditional medicine
Although practitioners of traditional medicine have used elderberry over centuries, there is little high-quality clinical evidence that such practices provide benefits.
Pigments
The pigments are used as colorants in various products, and "elderberry juice color" is listed by the US Food and Drug Administration as allowable in certified organic food products. In Japan, elderberry juice is listed as an approved "natural color additive" under the Food and Sanitation Law. Fibers can be dyed with elderberry juice (using alum as a mordant) to give a "muted purple" shade.
Other
The berry of S. callicarpa can be made into wine. The flowers of S. nigra are used to produce elderflower cordial. St-Germain, a French liqueur, is made from elderflowers. Hallands Fläder, a Swedish akvavit, is flavoured with elderflowers.
Hollowed elderberry twigs have traditionally been used as spiles to tap maple trees for syrup. Additionally, they have been hollowed out and used as flutes, blowguns, and syringes. In addition, the elderberry twigs and fruit are employed in creating dyes for basketry. These stems are dyed a very deep black by soaking them in a wash made from the berry stems of the elderberry.
The pith of elder has been used by watchmakers for cleaning tools before intricate work.
In culture
Folklore related to elder trees is extensive and can vary according to region. In some traditions, the elder tree is thought to ward off evil and give protection from witches, while other beliefs say that witches often congregate under the plant, especially when it is full of fruit. If an elder tree was cut down, a spirit known as the Elder Mother would be released and take her revenge. The tree could only safely be cut while chanting a rhyme to the Elder Mother. Romani people believe burning elder wood brings bad luck.
A wand made from the branch of an elder tree plays a pivotal role in the final book of the Harry Potter series, which was almost named Harry Potter and the Elder Wand.
Explanatory notes
Citations
General and cited references
| Biology and health sciences | Berries | Plants |
197767 | https://en.wikipedia.org/wiki/Radioactive%20decay | Radioactive decay | Radioactive decay (also known as nuclear decay, radioactivity, radioactive disintegration, or nuclear disintegration) is the process by which an unstable atomic nucleus loses energy by radiation. A material containing unstable nuclei is considered radioactive. Three of the most common types of decay are alpha, beta, and gamma decay. The weak force is the mechanism that is responsible for beta decay, while the other two are governed by the electromagnetic and nuclear forces.
Radioactive decay is a random process at the level of single atoms. According to quantum theory, it is impossible to predict when a particular atom will decay, regardless of how long the atom has existed. However, for a significant number of identical atoms, the overall decay rate can be expressed as a decay constant or as a half-life. The half-lives of radioactive atoms have a huge range: from nearly instantaneous to far longer than the age of the universe.
The decaying nucleus is called the parent radionuclide (or parent radioisotope), and the process produces at least one daughter nuclide. Except for gamma decay or internal conversion from a nuclear excited state, the decay is a nuclear transmutation resulting in a daughter containing a different number of protons or neutrons (or both). When the number of protons changes, an atom of a different chemical element is created.
There are 28 naturally occurring chemical elements on Earth that are radioactive, consisting of 35 radionuclides (seven elements have two different radionuclides each) that date before the time of formation of the Solar System. These 35 are known as primordial radionuclides. Well-known examples are uranium and thorium, but also included are naturally occurring long-lived radioisotopes, such as potassium-40. Each of the heavy primordial radionuclides participates in one of the four decay chains.
History of discovery
Henri Poincaré laid the seeds for the discovery of radioactivity through his interest in and studies of X-rays, which significantly influenced physicist Henri Becquerel. Radioactivity was discovered in 1896 by Becquerel and independently by Marie Curie, while working with phosphorescent materials. These materials glow in the dark after exposure to light, and Becquerel suspected that the glow produced in cathode-ray tubes by X-rays might be associated with phosphorescence. He wrapped a photographic plate in black paper and placed various phosphorescent salts on it. All results were negative until he used uranium salts. The uranium salts caused a blackening of the plate in spite of the plate being wrapped in black paper. These radiations were given the name "Becquerel Rays".
It soon became clear that the blackening of the plate had nothing to do with phosphorescence, as the blackening was also produced by non-phosphorescent salts of uranium and by metallic uranium. It became clear from these experiments that there was a form of invisible radiation that could pass through paper and was causing the plate to react as if exposed to light.
At first, it seemed as though the new radiation was similar to the then recently discovered X-rays. Further research by Becquerel, Ernest Rutherford, Paul Villard, Pierre Curie, Marie Curie, and others showed that this form of radioactivity was significantly more complicated. Rutherford was the first to realize that all such elements decay in accordance with the same mathematical exponential formula. Rutherford and his student Frederick Soddy were the first to realize that many decay processes resulted in the transmutation of one element to another. Subsequently, the radioactive displacement law of Fajans and Soddy was formulated to describe the products of alpha and beta decay.
The early researchers also discovered that many other chemical elements, besides uranium, have radioactive isotopes. A systematic search for the total radioactivity in uranium ores also guided Pierre and Marie Curie to isolate two new elements: polonium and radium. Except for the radioactivity of radium, the chemical similarity of radium to barium made these two elements difficult to distinguish.
Marie and Pierre Curie's study of radioactivity is an important factor in science and medicine. After their research on Becquerel's rays led them to the discovery of both radium and polonium, they coined the term "radioactivity" to define the emission of ionizing radiation by some heavy elements. (Later the term was generalized to all elements.) Their research on the penetrating rays in uranium and the discovery of radium launched an era of using radium for the treatment of cancer. Their exploration of radium could be seen as the first peaceful use of nuclear energy and the start of modern nuclear medicine.
Early health dangers
The dangers of ionizing radiation due to radioactivity and X-rays were not immediately recognized.
X-rays
The discovery of X‑rays by Wilhelm Röntgen in 1895 led to widespread experimentation by scientists, physicians, and inventors. Many people began recounting stories of burns, hair loss and worse in technical journals as early as 1896. In February of that year, Professor Daniel and Dr. Dudley of Vanderbilt University performed an experiment involving X-raying Dudley's head that resulted in his hair loss. A report by Dr. H.D. Hawks, of his suffering severe hand and chest burns in an X-ray demonstration, was the first of many other reports in Electrical Review.
Other experimenters, including Elihu Thomson and Nikola Tesla, also reported burns. Thomson deliberately exposed a finger to an X-ray tube over a period of time and suffered pain, swelling, and blistering. Other effects, including ultraviolet rays and ozone, were sometimes blamed for the damage, and many physicians still claimed that there were no effects from X-ray exposure at all.
Despite this, there were some early systematic hazard investigations, and as early as 1902 William Herbert Rollins wrote almost despairingly that his warnings about the dangers involved in the careless use of X-rays were not being heeded, either by industry or by his colleagues. By this time, Rollins had proved that X-rays could kill experimental animals, could cause a pregnant guinea pig to abort, and that they could kill a foetus. He also stressed that "animals vary in susceptibility to the external action of X-light" and warned that these differences be considered when patients were treated by means of X-rays.
Radioactive substances
However, the biological effects of radiation due to radioactive substances were less easy to gauge. This gave the opportunity for many physicians and corporations to market radioactive substances as patent medicines. Examples were radium enema treatments, and radium-containing waters to be drunk as tonics. Marie Curie protested against this sort of treatment, warning that "radium is dangerous in untrained hands". Curie later died from aplastic anaemia, likely caused by exposure to ionizing radiation. By the 1930s, after a number of cases of bone necrosis and death of radium treatment enthusiasts, radium-containing medicinal products had been largely removed from the market (radioactive quackery).
Radiation protection
Only a year after Röntgen's discovery of X-rays, the American engineer Wolfram Fuchs (1896) gave what is probably the first protection advice, but it was not until 1925 that the first International Congress of Radiology (ICR) was held and considered establishing international protection standards. The effects of radiation on genes, including the effect of cancer risk, were recognized much later. In 1927, Hermann Joseph Muller published research showing genetic effects and, in 1946, was awarded the Nobel Prize in Physiology or Medicine for his findings.
The second ICR was held in Stockholm in 1928 and proposed the adoption of the röntgen unit, and the International X-ray and Radium Protection Committee (IXRPC) was formed. Rolf Sievert was named chairman, but a driving force was George Kaye of the British National Physical Laboratory. The committee met in 1931, 1934, and 1937.
After World War II, the increased range and quantity of radioactive substances being handled as a result of military and civil nuclear programs led to large groups of occupational workers and the public being potentially exposed to harmful levels of ionising radiation. This was considered at the first post-war ICR convened in London in 1950, when the present International Commission on Radiological Protection (ICRP) was born.
Since then the ICRP has developed the present international system of radiation protection, covering all aspects of radiation hazards.
In 2020, Hauptmann and another 15 international researchers from eight nations (among them: Institutes of Biostatistics, Registry Research, Centers of Cancer Epidemiology, Radiation Epidemiology, and also the U.S. National Cancer Institute (NCI), International Agency for Research on Cancer (IARC) and the Radiation Effects Research Foundation of Hiroshima) studied definitively through meta-analysis the damage resulting from the "low doses" that have afflicted survivors of the atomic bombings of Hiroshima and Nagasaki and also in numerous accidents at nuclear plants that have occurred. These scientists reported, in JNCI Monographs: Epidemiological Studies of Low Dose Ionizing Radiation and Cancer Risk, that the new epidemiological studies directly support excess cancer risks from low-dose ionizing radiation. In 2021, Italian researcher Sebastiano Venturi reported the first correlations between radio-caesium and pancreatic cancer with the role of caesium in biology, in pancreatitis and in diabetes of pancreatic origin.
Units
The International System of Units (SI) unit of radioactive activity is the becquerel (Bq), named in honor of the scientist Henri Becquerel. One Bq is defined as one transformation (or decay or disintegration) per second.
An older unit of radioactivity is the curie, Ci, which was originally defined as "the quantity or mass of radium emanation in equilibrium with one gram of radium (element)". Today, the curie is defined as disintegrations per second, so that 1 curie (Ci) = .
For radiological protection purposes, although the United States Nuclear Regulatory Commission permits the use of the unit curie alongside SI units, the European Union European units of measurement directives required that its use for "public health ... purposes" be phased out by 31 December 1985.
The effects of ionizing radiation are often measured in units of gray for mechanical or sievert for damage to tissue.
Types
Radioactive decay results in a reduction of summed rest mass, once the released energy (the disintegration energy) has escaped in some way. Although decay energy is sometimes defined as associated with the difference between the mass of the parent nuclide products and the mass of the decay products, this is true only of rest mass measurements, where some energy has been removed from the product system. This is true because the decay energy must always carry mass with it, wherever it appears (see mass in special relativity) according to the formula E = mc2. The decay energy is initially released as the energy of emitted photons plus the kinetic energy of massive emitted particles (that is, particles that have rest mass). If these particles come to thermal equilibrium with their surroundings and photons are absorbed, then the decay energy is transformed to thermal energy, which retains its mass.
Decay energy, therefore, remains associated with a certain measure of the mass of the decay system, called invariant mass, which does not change during the decay, even though the energy of decay is distributed among decay particles. The energy of photons, the kinetic energy of emitted particles, and, later, the thermal energy of the surrounding matter, all contribute to the invariant mass of the system. Thus, while the sum of the rest masses of the particles is not conserved in radioactive decay, the system mass and system invariant mass (and also the system total energy) is conserved throughout any decay process. This is a restatement of the equivalent laws of conservation of energy and conservation of mass.
Alpha, beta and gamma decay
Early researchers found that an electric or magnetic field could split radioactive emissions into three types of beams. The rays were given the names alpha, beta, and gamma, in increasing order of their ability to penetrate matter. Alpha decay is observed only in heavier elements of atomic number 52 (tellurium) and greater, with the exception of beryllium-8 (which decays to two alpha particles). The other two types of decay are observed in all the elements. Lead, atomic number 82, is the heaviest element to have any isotopes stable (to the limit of measurement) to radioactive decay. Radioactive decay is seen in all isotopes of all elements of atomic number 83 (bismuth) or greater. Bismuth-209, however, is only very slightly radioactive, with a half-life greater than the age of the universe; radioisotopes with extremely long half-lives are considered effectively stable for practical purposes.
In analyzing the nature of the decay products, it was obvious from the direction of the electromagnetic forces applied to the radiations by external magnetic and electric fields that alpha particles carried a positive charge, beta particles carried a negative charge, and gamma rays were neutral. From the magnitude of deflection, it was clear that alpha particles were much more massive than beta particles. Passing alpha particles through a very thin glass window and trapping them in a discharge tube allowed researchers to study the emission spectrum of the captured particles, and ultimately proved that alpha particles are helium nuclei. Other experiments showed beta radiation, resulting from decay and cathode rays, were high-speed electrons. Likewise, gamma radiation and X-rays were found to be high-energy electromagnetic radiation.
The relationship between the types of decays also began to be examined: For example, gamma decay was almost always found to be associated with other types of decay, and occurred at about the same time, or afterwards. Gamma decay as a separate phenomenon, with its own half-life (now termed isomeric transition), was found in natural radioactivity to be a result of the gamma decay of excited metastable nuclear isomers, which were in turn created from other types of decay. Although alpha, beta, and gamma radiations were most commonly found, other types of emission were eventually discovered. Shortly after the discovery of the positron in cosmic ray products, it was realized that the same process that operates in classical beta decay can also produce positrons (positron emission), along with neutrinos (classical beta decay produces antineutrinos).
Electron capture
In electron capture, some proton-rich nuclides were found to capture their own atomic electrons instead of emitting positrons, and subsequently, these nuclides emit only a neutrino and a gamma ray from the excited nucleus (and often also Auger electrons and characteristic X-rays, as a result of the re-ordering of electrons to fill the place of the missing captured electron). These types of decay involve the nuclear capture of electrons or emission of electrons or positrons, and thus acts to move a nucleus toward the ratio of neutrons to protons that has the least energy for a given total number of nucleons. This consequently produces a more stable (lower energy) nucleus.
A hypothetical process of positron capture, analogous to electron capture, is theoretically possible in antimatter atoms, but has not been observed, as complex antimatter atoms beyond antihelium are not experimentally available. Such a decay would require antimatter atoms at least as complex as beryllium-7, which is the lightest known isotope of normal matter to undergo decay by electron capture.
Nucleon emission
Shortly after the discovery of the neutron in 1932, Enrico Fermi realized that certain rare beta-decay reactions immediately yield neutrons as an additional decay particle, so called beta-delayed neutron emission. Neutron emission usually happens from nuclei that are in an excited state, such as the excited 17O* produced from the beta decay of 17N. The neutron emission process itself is controlled by the nuclear force and therefore is extremely fast, sometimes referred to as "nearly instantaneous". Isolated proton emission was eventually observed in some elements. It was also found that some heavy elements may undergo spontaneous fission into products that vary in composition. In a phenomenon called cluster decay, specific combinations of neutrons and protons other than alpha particles (helium nuclei) were found to be spontaneously emitted from atoms.
More exotic types of decay
Other types of radioactive decay were found to emit previously seen particles but via different mechanisms. An example is internal conversion, which results in an initial electron emission, and then often further characteristic X-rays and Auger electrons emissions, although the internal conversion process involves neither beta nor gamma decay. A neutrino is not emitted, and none of the electron(s) and photon(s) emitted originate in the nucleus, even though the energy to emit all of them does originate there. Internal conversion decay, like isomeric transition gamma decay and neutron emission, involves the release of energy by an excited nuclide, without the transmutation of one element into another.
Rare events that involve a combination of two beta-decay-type events happening simultaneously are known (see below). Any decay process that does not violate the conservation of energy or momentum laws (and perhaps other particle conservation laws) is permitted to happen, although not all have been detected. An interesting example discussed in a final section, is bound state beta decay of rhenium-187. In this process, the beta electron-decay of the parent nuclide is not accompanied by beta electron emission, because the beta particle has been captured into the K-shell of the emitting atom. An antineutrino is emitted, as in all negative beta decays.
If energy circumstances are favorable, a given radionuclide may undergo many competing types of decay, with some atoms decaying by one route, and others decaying by another. An example is copper-64, which has 29 protons, and 35 neutrons, which decays with a half-life of hours. This isotope has one unpaired proton and one unpaired neutron, so either the proton or the neutron can decay to the other particle, which has opposite isospin. This particular nuclide (though not all nuclides in this situation) is more likely to decay through beta plus decay (%) than through electron capture (%). The excited energy states resulting from these decays which fail to end in a ground energy state, also produce later internal conversion and gamma decay in almost 0.5% of the time.
List of decay modes
Decay chains and multiple modes
The daughter nuclide of a decay event may also be unstable (radioactive). In this case, it too will decay, producing radiation. The resulting second daughter nuclide may also be radioactive. This can lead to a sequence of several decay events called a decay chain (see this article for specific details of important natural decay chains). Eventually, a stable nuclide is produced. Any decay daughters that are the result of an alpha decay will also result in helium atoms being created.
Some radionuclides may have several different paths of decay. For example, % of bismuth-212 decays, through alpha-emission, to thallium-208 while % of bismuth-212 decays, through beta-emission, to polonium-212. Both thallium-208 and polonium-212 are radioactive daughter products of bismuth-212, and both decay directly to stable lead-208.
Occurrence and applications
According to the Big Bang theory, stable isotopes of the lightest three elements (H, He, and traces of Li) were produced very shortly after the emergence of the universe, in a process called Big Bang nucleosynthesis. These lightest stable nuclides (including deuterium) survive to today, but any radioactive isotopes of the light elements produced in the Big Bang (such as tritium) have long since decayed. Isotopes of elements heavier than boron were not produced at all in the Big Bang, and these first five elements do not have any long-lived radioisotopes. Thus, all radioactive nuclei are, therefore, relatively young with respect to the birth of the universe, having formed later in various other types of nucleosynthesis in stars (in particular, supernovae), and also during ongoing interactions between stable isotopes and energetic particles. For example, carbon-14, a radioactive nuclide with a half-life of only years, is constantly produced in Earth's upper atmosphere due to interactions between cosmic rays and nitrogen.
Nuclides that are produced by radioactive decay are called radiogenic nuclides, whether they themselves are stable or not. There exist stable radiogenic nuclides that were formed from short-lived extinct radionuclides in the early Solar System. The extra presence of these stable radiogenic nuclides (such as xenon-129 from extinct iodine-129) against the background of primordial stable nuclides can be inferred by various means.
Radioactive decay has been put to use in the technique of radioisotopic labeling, which is used to track the passage of a chemical substance through a complex system (such as a living organism). A sample of the substance is synthesized with a high concentration of unstable atoms. The presence of the substance in one or another part of the system is determined by detecting the locations of decay events.
On the premise that radioactive decay is truly random (rather than merely chaotic), it has been used in hardware random-number generators. Because the process is not thought to vary significantly in mechanism over time, it is also a valuable tool in estimating the absolute ages of certain materials. For geological materials, the radioisotopes and some of their decay products become trapped when a rock solidifies, and can then later be used (subject to many well-known qualifications) to estimate the date of the solidification. These include checking the results of several simultaneous processes and their products against each other, within the same sample. In a similar fashion, and also subject to qualification, the rate of formation of carbon-14 in various eras, the date of formation of organic matter within a certain period related to the isotope's half-life may be estimated, because the carbon-14 becomes trapped when the organic matter grows and incorporates the new carbon-14 from the air. Thereafter, the amount of carbon-14 in organic matter decreases according to decay processes that may also be independently cross-checked by other means (such as checking the carbon-14 in individual tree rings, for example).
Szilard–Chalmers effect
The Szilard–Chalmers effect is the breaking of a chemical bond as a result of a kinetic energy imparted from radioactive decay. It operates by the absorption of neutrons by an atom and subsequent emission of gamma rays, often with significant amounts of kinetic energy. This kinetic energy, by Newton's third law, pushes back on the decaying atom, which causes it to move with enough speed to break a chemical bond. This effect can be used to separate isotopes by chemical means.
The Szilard–Chalmers effect was discovered in 1934 by Leó Szilárd and Thomas A. Chalmers. They observed that after bombardment by neutrons, the breaking of a bond in liquid ethyl iodide allowed radioactive iodine to be removed.
Origins of radioactive nuclides
Radioactive primordial nuclides found in the Earth are residues from ancient supernova explosions that occurred before the formation of the Solar System. They are the fraction of radionuclides that survived from that time, through the formation of the primordial solar nebula, through planet accretion, and up to the present time. The naturally occurring short-lived radiogenic radionuclides found in today's rocks, are the daughters of those radioactive primordial nuclides. Another minor source of naturally occurring radioactive nuclides are cosmogenic nuclides, that are formed by cosmic ray bombardment of material in the Earth's atmosphere or crust. The decay of the radionuclides in rocks of the Earth's mantle and crust contribute significantly to Earth's internal heat budget.
Aggregate processes
While the underlying process of radioactive decay is subatomic, historically and in most practical cases it is encountered in bulk materials with very large numbers of atoms. This section discusses models that connect events at the atomic level to observations in aggregate.
Terminology
The decay rate, or activity, of a radioactive substance is characterized by the following time-independent parameters:
The half-life, , is the time taken for the activity of a given amount of a radioactive substance to decay to half of its initial value.
The decay constant, "lambda", the reciprocal of the mean lifetime (in ), sometimes referred to as simply decay rate.
The mean lifetime, "tau", the average lifetime (1/e life) of a radioactive particle before decay.
Although these are constants, they are associated with the statistical behavior of populations of atoms. In consequence, predictions using these constants are less accurate for minuscule samples of atoms.
In principle a half-life, a third-life, or even a (1/√2)-life, could be used in exactly the same way as half-life; but the mean life and half-life have been adopted as standard times associated with exponential decay.
Those parameters can be related to the following time-dependent parameters:
Total activity (or just activity), , is the number of decays per unit time of a radioactive sample.
Number of particles, , in the sample.
Specific activity, , is the number of decays per unit time per amount of substance of the sample at time set to zero (). "Amount of substance" can be the mass, volume or moles of the initial sample.
These are related as follows:
where N0 is the initial amount of active substance — substance that has the same percentage of unstable particles as when the substance was formed.
Assumptions
The mathematics of radioactive decay depend on a key assumption that a nucleus of a radionuclide has no "memory" or way of translating its history into its present behavior. A nucleus does not "age" with the passage of time. Thus, the probability of its breaking down does not increase with time but stays constant, no matter how long the nucleus has existed. This constant probability may differ greatly between one type of nucleus and another, leading to the many different observed decay rates. However, whatever the probability is, it does not change over time. This is in marked contrast to complex objects that do show aging, such as automobiles and humans. These aging systems do have a chance of breakdown per unit of time that increases from the moment they begin their existence.
Aggregate processes, like the radioactive decay of a lump of atoms, for which the single-event probability of realization is very small but in which the number of time-slices is so large that there is nevertheless a reasonable rate of events, are modelled by the Poisson distribution, which is discrete. Radioactive decay and nuclear particle reactions are two examples of such aggregate processes. The mathematics of Poisson processes reduce to the law of exponential decay, which describes the statistical behaviour of a large number of nuclei, rather than one individual nucleus. In the following formalism, the number of nuclei or the nuclei population N, is of course a discrete variable (a natural number)—but for any physical sample N is so large that it can be treated as a continuous variable. Differential calculus is used to model the behaviour of nuclear decay.
One-decay process
Consider the case of a nuclide that decays into another by some process (emission of other particles, like electron neutrinos and electrons e− as in beta decay, are irrelevant in what follows). The decay of an unstable nucleus is entirely random in time so it is impossible to predict when a particular atom will decay. However, it is equally likely to decay at any instant in time. Therefore, given a sample of a particular radioisotope, the number of decay events expected to occur in a small interval of time is proportional to the number of atoms present , that is
Particular radionuclides decay at different rates, so each has its own decay constant . The expected decay is proportional to an increment of time, :
The negative sign indicates that decreases as time increases, as the decay events follow one after another. The solution to this first-order differential equation is the function:
where is the value of at time = 0, with the decay constant expressed as
We have for all time :
where is the constant number of particles throughout the decay process, which is equal to the initial number of nuclides since this is the initial substance.
If the number of non-decayed nuclei is:
then the number of nuclei of (i.e. the number of decayed nuclei) is
The number of decays observed over a given interval obeys Poisson statistics. If the average number of decays is , the probability of a given number of decays is
Chain-decay processes
Chain of two decays
Now consider the case of a chain of two decays: one nuclide decaying into another by one process, then decaying into another by a second process, i.e. . The previous equation cannot be applied to the decay chain, but can be generalized as follows. Since decays into , then decays into , the activity of adds to the total number of nuclides in the present sample, before those nuclides decay and reduce the number of nuclides leading to the later sample. In other words, the number of second generation nuclei increases as a result of the first generation nuclei decay of , and decreases as a result of its own decay into the third generation nuclei . The sum of these two terms gives the law for a decay chain for two nuclides:
The rate of change of , that is , is related to the changes in the amounts of and , can increase as is produced from and decrease as produces .
Re-writing using the previous results:
The subscripts simply refer to the respective nuclides, i.e. is the number of nuclides of type ; is the initial number of nuclides of type ; is the decay constant for – and similarly for nuclide . Solving this equation for gives:
In the case where is a stable nuclide ( = 0), this equation reduces to the previous solution:
as shown above for one decay. The solution can be found by the integration factor method, where the integrating factor is . This case is perhaps the most useful since it can derive both the one-decay equation (above) and the equation for multi-decay chains (below) more directly.
Chain of any number of decays
For the general case of any number of consecutive decays in a decay chain, i.e. , where is the number of decays and is a dummy index (), each nuclide population can be found in terms of the previous population. In this case , , ..., . Using the above result in a recursive form:
The general solution to the recursive problem is given by Bateman's equations:
Multiple products
In all of the above examples, the initial nuclide decays into just one product. Consider the case of one initial nuclide that can decay into either of two products, that is and in parallel. For example, in a sample of potassium-40, 89.3% of the nuclei decay to calcium-40 and 10.7% to argon-40. We have for all time :
which is constant, since the total number of nuclides remains constant. Differentiating with respect to time:
defining the total decay constant in terms of the sum of partial decay constants and :
Solving this equation for :
where is the initial number of nuclide A. When measuring the production of one nuclide, one can only observe the total decay constant . The decay constants and determine the probability for the decay to result in products or as follows:
because the fraction of nuclei decay into while the fraction of nuclei decay into .
Corollaries of laws
The above equations can also be written using quantities related to the number of nuclide particles in a sample;
The activity: .
The amount of substance: .
The mass: .
where = is the Avogadro constant, is the molar mass of the substance in kg/mol, and the amount of the substance is in moles.
Decay timing: definitions and relations
Time constant and mean-life
For the one-decay solution :
the equation indicates that the decay constant has units of , and can thus also be represented as 1/, where is a characteristic time of the process called the time constant.
In a radioactive decay process, this time constant is also the mean lifetime for decaying atoms. Each atom "lives" for a finite amount of time before it decays, and it may be shown that this mean lifetime is the arithmetic mean of all the atoms' lifetimes, and that it is , which again is related to the decay constant as follows:
This form is also true for two-decay processes simultaneously , inserting the equivalent values of decay constants (as given above)
into the decay solution leads to:
Half-life
A more commonly used parameter is the half-life . Given a sample of a particular radionuclide, the half-life is the time taken for half the radionuclide's atoms to decay. For the case of one-decay nuclear reactions:
the half-life is related to the decay constant as follows: set and = to obtain
This relationship between the half-life and the decay constant shows that highly radioactive substances are quickly spent, while those that radiate weakly endure longer. Half-lives of known radionuclides vary by almost 54 orders of magnitude, from more than years ( sec) for the very nearly stable nuclide 128Te, to seconds for the highly unstable nuclide 5H.
The factor of in the above relations results from the fact that the concept of "half-life" is merely a way of selecting a different base other than the natural base for the lifetime expression. The time constant is the -life, the time until only 1/e remains, about 36.8%, rather than the 50% in the half-life of a radionuclide. Thus, is longer than . The following equation can be shown to be valid:
Since radioactive decay is exponential with a constant probability, each process could as easily be described with a different constant time period that (for example) gave its "(1/3)-life" (how long until only 1/3 is left) or "(1/10)-life" (a time period until only 10% is left), and so on. Thus, the choice of and for marker-times, are only for convenience, and from convention. They reflect a fundamental principle only in so much as they show that the same proportion of a given radioactive substance will decay, during any time-period that one chooses.
Mathematically, the life for the above situation would be found in the same way as aboveby setting , and substituting into the decay solution to obtain
Example for carbon-14
Carbon-14 has a half-life of years and a decay rate of 14 disintegrations per minute (dpm) per gram of natural carbon.
If an artifact is found to have radioactivity of 4 dpm per gram of its present C, we can find the approximate age of the object using the above equation:
where:
Changing rates
The radioactive decay modes of electron capture and internal conversion are known to be slightly sensitive to chemical and environmental effects that change the electronic structure of the atom, which in turn affects the presence of 1s and 2s electrons that participate in the decay process. A small number of nuclides are affected. For example, chemical bonds can affect the rate of electron capture to a small degree (in general, less than 1%) depending on the proximity of electrons to the nucleus. In 7Be, a difference of 0.9% has been observed between half-lives in metallic and insulating environments. This relatively large effect is because beryllium is a small atom whose valence electrons are in 2s atomic orbitals, which are subject to electron capture in 7Be because (like all s atomic orbitals in all atoms) they naturally penetrate into the nucleus.
In 1992, Jung et al. of the Darmstadt Heavy-Ion Research group observed an accelerated β− decay of 163Dy66+. Although neutral 163Dy is a stable isotope, the fully ionized 163Dy66+ undergoes β− decay into the K and L shells to 163Ho66+ with a half-life of 47 days.
Rhenium-187 is another spectacular example. 187Re normally undergoes beta decay to 187Os with a half-life of 41.6 × 109 years, but studies using fully ionised 187Re atoms (bare nuclei) have found that this can decrease to only 32.9 years. This is attributed to "bound-state β− decay" of the fully ionised atom – the electron is emitted into the "K-shell" (1s atomic orbital), which cannot occur for neutral atoms in which all low-lying bound states are occupied.
A number of experiments have found that decay rates of other modes of artificial and naturally occurring radioisotopes are, to a high degree of precision, unaffected by external conditions such as temperature, pressure, the chemical environment, and electric, magnetic, or gravitational fields. Comparison of laboratory experiments over the last century, studies of the Oklo natural nuclear reactor (which exemplified the effects of thermal neutrons on nuclear decay), and astrophysical observations of the luminosity decays of distant supernovae (which occurred far away so the light has taken a great deal of time to reach us), for example, strongly indicate that unperturbed decay rates have been constant (at least to within the limitations of small experimental errors) as a function of time as well.
Recent results suggest the possibility that decay rates might have a weak dependence on environmental factors. It has been suggested that measurements of decay rates of silicon-32, manganese-54, and radium-226 exhibit small seasonal variations (of the order of 0.1%). However, such measurements are highly susceptible to systematic errors, and a subsequent paper has found no evidence for such correlations in seven other isotopes (22Na, 44Ti, 108Ag, 121Sn, 133Ba, 241Am, 238Pu), and sets upper limits on the size of any such effects. The decay of radon-222 was once reported to exhibit large 4% peak-to-peak seasonal variations (see plot), which were proposed to be related to either solar flare activity or the distance from the Sun, but detailed analysis of the experiment's design flaws, along with comparisons to other, much more stringent and systematically controlled, experiments refute this claim.
GSI anomaly
An unexpected series of experimental results for the rate of decay of heavy highly charged radioactive ions circulating in a storage ring has provoked theoretical activity in an effort to find a convincing explanation. The rates of weak decay of two radioactive species with half lives of about 40 s and 200 s are found to have a significant oscillatory modulation, with a period of about 7 s.
The observed phenomenon is known as the GSI anomaly, as the storage ring is a facility at the GSI Helmholtz Centre for Heavy Ion Research in Darmstadt, Germany. As the decay process produces an electron neutrino, some of the proposed explanations for the observed rate oscillation invoke neutrino properties. Initial ideas related to flavour oscillation met with skepticism. A more recent proposal involves mass differences between neutrino mass eigenstates.
Nuclear processes
A nuclide is considered to "exist" if it has a half-life greater than 2x10−14s. This is an arbitrary boundary; shorter half-lives are considered resonances, such as a system undergoing a nuclear reaction. This time scale is characteristic of the strong interaction which creates the nuclear force. Only nuclides are considered to decay and produce radioactivity.
Nuclides can be stable or unstable. Unstable nuclides decay, possibly in several steps, until they become stable. There are 251 known stable nuclides. The number of unstable nuclides discovered has grown, with about 3000 known in 2006.
The most common and consequently historically the most important forms of natural radioactive decay involve the emission of alpha-particles, beta-particles, and gamma rays. Each of these correspond to a fundamental interaction predominantly responsible for the radioactivity:
alpha-decay -> strong interaction,
beta-decay -> weak interaction,
gamma-decay -> electromagnetism.
In alpha decay, a particle containing two protons and two neutrons, equivalent to a He nucleus, breaks out of the parent nucleus. The process represents a competition between the electromagnetic repulsion between the protons in the nucleus and attractive nuclear force, a residual of the strong interaction. The alpha particle is an especially strongly bound nucleus, helping it win the competition more often. However some nuclei break up or fission into larger particles and artificial nuclei decay with the emission of
single protons, double protons, and other combinations.
Beta decay transforms a neutron into proton or vice versa. When a neutron inside a parent nuclide decays to a proton, an electron, a anti-neutrino, and nuclide with high atomic number results. When a proton in a parent nuclide transforms to a neutron, a positron, a neutrino, and nuclide with a lower atomic number results. These changes are a direct manifestation of the weak interaction.
Gamma decay resembles other kinds of electromagnetic emission: it corresponds to transitions between an excited quantum state and lower energy state. Any of the particle decay mechanisms often leave the daughter in an excited state, which then decays via gamma emission.
Other forms of decay include neutron emission, electron capture, internal conversion, cluster decay.
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197817 | https://en.wikipedia.org/wiki/General%20practitioner | General practitioner | A general practitioner (GP) or family physician is a doctor who is a consultant in general practice.
GPs have distinct expertise and experience in providing whole person medical care, whilst managing the complexity, uncertainty and risk associated with the continuous care they provide. GPs work at the heart of their communities, striving to provide comprehensive and equitable care for everyone, taking into account their health care needs, stage of life and background. GPs work in, connect with and lead multidisciplinary teams that care for people and their families, respecting the context in which they live, aiming to ensure all of their physical health and mental health needs are met. They are trained to treat patients to levels of complexity that vary between countries. The term "primary care physician" is used in the United States.
A core element in general practice is continuity of care, that bridges episodes of various illnesses over time. Greater continuity with a general practitioner has been shown to reduce the need for out-of-hours services and acute hospital admittance. Continuous care by the same general practitioner has been found to reduce mortality.
The role of a GP varies between and within countries, and is often dependent on local needs and circumstances. In urban areas their roles may focus on:
care of chronic/complex health conditions
treatment of urgent/acute non-life-threatening diseases
mental health care
preventive care, including health education and immunisation.
screening/early detection of disease
palliative care
care coordination/referral to allied health professions or specialised medical care
In rural areas, a GP may additionally be routinely involved in pre-hospital emergency care, the delivery of babies, community hospital care and performing low-complexity surgical procedures. GPs may work in larger primary care centers where they provide care within a multidisciplinary healthcare team, while in other cases GPs may work as sole practitioners or in smaller practices.
The term general practitioner or GP is common in the United Kingdom, Republic of Ireland, Australia, Canada, Singapore, South Africa, New Zealand and other Commonwealth countries. In these countries, the word "physician" is largely reserved for medical specialists often working in hospitals, notably in internal medicine. In North America the term is sometimes synonymous with the terms family doctor or primary care physician.
General practice is an academic and scientific discipline with its own educational content, research, evidence base and clinical activity. Historically, the role of a GP was performed by any doctor with qualifications from a medical school working in the community. However, since the 1950s, general practice has become a medical specialty with additional training requirements. The 1978 Alma Ata Declaration set the intellectual foundation of primary care and general practice.
Asia
India
The basic medical degree in India is MBBS (Bachelor of Medicine, Bachelor of Surgery). These generally consist of a four-and-a-half-year course followed by a year of compulsory rotatory internship in India. The internship requires the candidate to work in all departments for a stipulated period of time, to undergo hands-on training in treating patients.
The registration of doctors is usually managed by state medical councils. A permanent registration as a Registered Medical Practitioner is granted only after satisfactory completion of the compulsory internship.
The Federation of Family Physicians' Associations of India (FFPAI) is an organization which has a connection with more than 8000 general practitioners through having affiliated membership.
Bangladesh
In Bangladesh, the completion of a 5-year MBBS program is succeeded by a one-year rotational internship encompassing various specialties. Bangladesh Medical & Dental Council (BM&DC) then provides permanent registration to the doctors, after which the candidate may choose to practice as a GP or opt for specialty training.. As of 2019, there are some 86,800 doctors, and dentists registered with the BM&DC.
Bangladesh College of Physicians and Surgeons (BCPS) has a one-year membership and four-year fellowship program in Family Medicine.
Pakistan
In Pakistan, 5 years of MBBS is followed by one year of internship in different specialties. Pakistan Medical and Dental Council (PMDC) then confers permanent registration, after which the candidate may choose to practice as a GP or opt for specialty training.
The first Family Medicine Training programme was approved by the College of Physicians and Surgeons of Pakistan (CPSP) in 1992 and initiated in 1993 by the Family Medicine Division of the Department of Community Health
Sciences, Aga Khan University, Pakistan.
Family Medicine residency training programme of Ziauddin University is approved for Fellowship in Family Medicine.
Europe
European Union
General practitioners are regulated in the EU by Directive 2005/36/EC.
France
In France the médecin généraliste (commonly called médecin de famille) is responsible for primary care medicine, including non-vital emergencies and patient's follow up. This implies prevention, education, care of the diseases and traumas. The general practitionner orientates the patients to other specialists when necessary.
They have a role in the survey of epidemics, a legal role (constatation of traumas that can bring compensation, certificates for the practice of a sport, death certificate, certificate for hospitalisation without consent in case of mental incapacity), and a role in the emergency care (they can be called by the samu, the French EMS). They often go to a patient's home when the patient cannot come to the consulting room (especially in case of children or old people), and have to contribute to night and week-end duties.
The studies consist of six years at university (common to all medical specialities), and four years as a resident (interne) :
the first year (PASS, Parcours d'Accès Spécifique Santé, often abbreviated to P1 by students) is common with the dentists, pharmacists and midwives. The rank at the final competitive examination determines in which branch the student can choose to study.
the following two years, called propédeutique, are dedicated to the fundamental sciences: anatomy, human physiology, biochemistry, bacteriology, statistics...
the three following years are called externat and are dedicated to the study of clinical medicine; they end with a classifying examination, the rank determines in which specialty (general medicine is one of them) the student can make her or his internat;
the internat is three years -or more depending on the specialty- of initial professional experience under the responsibility of a senior; the interne can prescribe, s/he can replace physicians, and usually works in a hospital.
This ends with a doctorate, a research work which usually consist of a statistical study of cases to propose a care strategy for a specific condition (in an epidemiological, diagnostic, or therapeutic point of view).
Greece
General Practice was established as a medical specialty in Greece in 1986. To qualify as a General Practitioner (γενικός ιατρός, genikos iatros) doctors in Greece are required to complete four years of vocational training after medical school, including three years and two months in a hospital setting. General Practitioners in Greece may either work as private specialists or for the National Healthcare Service, ESY (Εθνικό Σύστημα Υγείας, ΕΣΥ).
Netherlands and Belgium
General practice in the Netherlands and Belgium is considered advanced. The huisarts (literally: "home doctor") administers first line, primary care. In the Netherlands, patients usually cannot consult a hospital specialist without a required referral. Most GPs work in private practice although more medical centers with employed GPs are seen. Many GPs have a specialist interest, e.g. in palliative care.
In Belgium, one year of lectures and two years of residency are required. In the Netherlands, training consists of three years (full-time) of specialization after completion of internships of 3 years. First and third year of training takes place at a GP practice. The second year of training consists of six months training at an emergency room, or internal medicine, paediatrics or gynaecology, or a combination of a general or academic hospital, three months of training at a psychiatric hospital or outpatient clinic and three months at a nursing home (verpleeghuis) or clinical geriatrics ward/policlinic. During all three years, residents get one day of training at university while working in practice the other days. The first year, a lot of emphasis is placed on communications skills with video training. Furthermore, all aspects of working as a GP gets addressed including working with the medical standards from the Dutch GP association NHG (Nederlands Huisartsen Genootschap).
All residents must also take the national GP knowledge test (Landelijke Huisartsgeneeskundige Kennistoets (LHK-toets)) twice a year. In this test of 120 multiple choice questions, medical, ethical, scientific and legal matters of GP work are addressed.
Spain
In Spain GPs are officially especialistas en medicina familiar y comunitaria but are commonly called "médico de cabecera" or "médico de familia". It was established as a medical specialty in Spain in 1978.
Most Spanish GPs work for the state-funded health services provided by the county's 17 regional governments (comunidades autónomas). They are in most cases salary-based healthcare workers.
For the provision of primary care, Spain is currently divided geographically in basic health care areas (áreas básicas de salud), each one containing a primary health care team (Equipo de atención primaria). Each team is multidisciplinary and typically includes GPs, community pediatricians, nurses, physiotherapists and social workers, together with ancillary staff. In urban areas all the services are concentrated in a single large building (Centro de salud) while in rural areas the main center is supported by smaller branches (consultorios), typically single-handled.
Becoming a GP in Spain involves studying medicine for 6 years, passing a competitive national exam called MIR (Medico Interno Residente) and undergoing a 4-year training program. The training program includes core specialties as general medicine and general practice (around 12 months each), pediatrics, gynecology, orthopedics and psychiatry. Shorter and optional placements in ENT, ophthalmology, ED, infectious diseases, rheumathology or others add up to the 4 years curriculum. The assessment is work based and involves completing a logbook that ensures all the expected skills, abilities and aptitudes have been acquired by the end of the training period.
Russia
In the Russian Federation, the General Practitioner's Regulation was put into effect in 1992, after which medical schools started training in the relevant specialty.
The right to practice as a general practitioner gives a certificate of appropriate qualifications. General medical practice can be carried out both individually and in a group, including with the participation of narrow specialists. The work of general practitioners is allowed, both in the medical institution and in private.
The general practitioner has broad legal rights. He can lead junior medical personnel, provide services under medical insurance contracts, conclude additional contracts to the main contract, and conduct an examination of the quality of medical services. For independent decisions, the general practitioner is responsible in accordance with the law.
The main tasks of a general practitioner are:
Prevention, diagnosis and treatment of the most common diseases;
Emergency and emergency medical care;
Performance of medical manipulations.
United Kingdom
In the United Kingdom physicians wishing to become GPs take at least five years' training after medical school, which is usually an undergraduate course of five to six years (or a graduate course of four to six years) leading to the degrees of Bachelor of Medicine and Bachelor of Surgery.
Until 2005, those wishing to become a general practitioner of medicine had to do a minimum of the following postgraduate training:
One year as a pre-registration house officer (PRHO) (formerly called a house officer), in which the trainee would usually spend six months on a general surgical ward and six months on a general medical ward in a hospital;
Two years as a senior house officer (SHO) – often on a General Practice Vocational Training Scheme (GP-VTS) in which the trainee would normally complete four six-month jobs in hospital specialties such as obstetrics and gynaecology, paediatrics, geriatric medicine, accident and emergency or psychiatry;
One year as a general practice registrar on a GPST.
This process changed under the programme Modernising Medical Careers. Medical practitioners graduating from 2005 onwards have to do a minimum of five years postgraduate training:
Two years of Foundation Training, in which the trainee will do a rotation around either six four-month jobs or eight three-month jobs – these include at least three months in general medicine and three months in general surgery, but will also include jobs in other areas;
A three-year "run-through" GP Speciality Training Programme containing (GPSTP): This comprises a minimum of twelve months as a hospital based Specialty Trainee during which time the trainee completes a mixture of jobs in specialties such as obstetrics and gynaecology, paediatrics, geriatric medicine, accident and emergency or psychiatry; eighteen to twenty-four months as a GP Specialty Trainee working in General Practice. The balance of training time spent in hospital versus in GP is planned to shift in 2022 to be consistently 12 months' hospital training and 24 months' training time in general practice.
The postgraduate qualification Membership of the Royal College of General Practitioners (MRCGP) was previously optional. In 2008, a requirement was introduced for physicians to succeed in the MRCGP assessments in order to be issued with a certificate of completion of their specialty training (CCT) in general practice. After passing the assessments, they are eligible to use the post-nominal letters MRCGP (so long as the physicians continued to pay membership fees to the RCGP, though many do not). During the GP specialty training programme, the medical practitioner must complete a variety of assessments in order to be allowed to practice independently as a GP. There is a knowledge-based exam with multiple choice questions called the Applied Knowledge Test (AKT). The practical examination takes the form of a "simulated surgery" in which the physicians is presented with thirteen clinical cases and assessment is made of data gathering, interpersonal skills and clinical management. This Clinical Skills Assessment (CSA) is held on three or four occasions throughout the year and takes place at the renovated headquarters of the Royal College of General Practitioners (RCGP), at 30 Euston Square, London. Finally, throughout the year the physician must complete an electronic portfolio which is made up of case-based discussions, critique of videoed consultations and reflective entries into a "learning log".
In addition, many hold qualifications such as the DCH (Diploma in Child Health of the Royal College of Paediatrics and Child Health) or the DRCOG (Diploma of the Royal College of Obstetricians and Gynaecologists), the DPD (Diploma in Practical Dermatology) or the DGH (Diploma in Geriatric Medicine of the Royal College of Physicians). Some General Practitioners also hold the MRCP (Member of the Royal College of Physicians) or other specialist qualifications, but generally only if they had a hospital career, or a career in another speciality, before training in General Practice.
There are many arrangements under which general practitioners can work in the UK. While the main career aim is becoming a principal or partner in a GP surgery, many become salaried or non-principal GPs, work in hospitals in GP-led acute care units or perform locum work. Whichever of these roles they fill, the vast majority of GPs receive most of their income from the National Health Service (NHS). Principals and partners in GP surgeries are self-employed, but they have contractual arrangements with the NHS which give them considerable predictability of income.
Visits to GP surgeries are free in all countries of the United Kingdom, but charges for prescriptions are applied in England (except for those over 60, under 18, and those on low incomes and welfare). Wales, Scotland and Northern Ireland have abolished all charges.
Recent reforms to the NHS have included changes to the GP contract. General practitioners are no longer required to work unsociable hours and get paid to some extent according to their performance, (e.g. numbers of patients treated, what treatments were administered, and the health of their catchment area, through the Quality and Outcomes Framework). The IT system used for assessing their income based on these criteria is called QMAS. The amount that a GP can expect to earn does vary according to the location of their work and the health needs of the population that they serve. Within a couple of years of the new contract being introduced, it became apparent that there were a few examples where the arrangements were out step with what had been expected. A full-time self-employed GP, such as a GMS or PMS practice partner, might currently expect to earn a profit share of around £95,900 before tax while a GP employed by a CCG could expect to earn a salary in the range of £54,863 to £82,789. This can equate to an hourly rate of around £40 an hour for a GP partner.
A survey by Ipsos MORI released in 2011 reports that 88% of adults in the UK "trust doctors to tell the truth".
In May 2017, there was said to be a crisis in the UK with practices having difficulties recruiting GPs they need. Helen Stokes-Lampard of the Royal College of General Practitioners said, "At present, UK general practice does not have sufficient resources to deliver the care and services necessary to meet our patients' changing needs, meaning that GPs and our teams are working under intense pressures, which are simply unsustainable. Workload in general practice is escalating – it has increased 16% over the last seven years, according to the latest research – yet investment in our service has steadily declined over the last decade and the number of GPs has not risen in step with patient demand ... This must be addressed as a matter of urgency.". Professor Azeem Majeed from Imperial College has also raised concerns about general practice in the UK.
In 2018 the average GP worked less than three and a half days a week because of the "intensity of working day".
There is an NHS England initiative to situate GPs in or near hospital emergency departments to divert minor cases away from A&E and reduce pressure on emergency services. 97 hospital trusts have been allocated money, mostly for premises alterations or development.
North America
United States
The population of this type of medical practitioner is declining, however. Currently, the Medical Departments of the US Air Force, Army and Navy have many of these general practitioners, known as General Medical Officers or GMOs, in active practice. The GMO is an inherent concept to all military medical branches. GMOs are the gatekeepers of medicine in that they hold the purse strings and decide upon the merit of specialist consultation. The US now holds a different definition for the term "general practitioner". The two terms "general practitioner" and "family practice" were synonymous prior to 1970. At that time both terms (if used within the US) referred to someone who completed medical school and the one-year required internship, and then worked as a general family doctor. Completion of a post-graduate specialty training program or residency in family medicine was, at that time, not a requirement. A physician who specializes in "family medicine" must now complete a residency in family medicine and must be eligible for board certification, which is required by many hospitals and health plans for hospital privileges and remuneration, respectively. It was not until the 1970s that family medicine was recognized as a specialty in the US.
Many licensed family medical practitioners in the United States after this change began to use the term "general practitioner" to refer to those practitioners who previously did not complete a family medicine residency. Family physicians (after completing medical school) must then complete three to four years of additional residency in family medicine. Three hundred hours of medical education within the prior six years is also required to be eligible to sit for the board certification exam; these hours are largely acquired during residency training.
The existing general practitioners in the 1970s were given the choice to be grandfathered into the newly created specialty of Family Practice. In 1971 the American Academy of General Practice changed its name to the American Academy of Family Physicians. The prior system of graduating from medical school and completing one year of post-graduate training (rotating internship) was not abolished as 47 of the 50 states allow a physician to obtain a medical license without completion of residency. If one wanted to become a "house-call-making" type of physician, one still needs to only complete one or two years of a residency in either pediatrics, family medicine or internal medicine. This would make a physician a non-board eligible general practitioner able to qualify and obtain a license to practice medicine in 47 of the 50 United States of America. Since the establishment of the Board of Family Medicine, a family medicine physician is no longer the same as a general practitioner. What makes a Family Medicine Physician different from a General Practitioner/Physician is two-fold. First off a Family Medicine Physician has completed the three years of Family Medicine residency and is board eligible or board certified in Family Medicine; while a General Practitioner does not have any board certification and cannot sit for any board exam. Secondly, a Family Medicine Physician is able to practice obstetrics, the care of the pregnant woman from conception to delivery, while a general practitioner is not adequately trained in obstetrics.
Prior to recent history most postgraduate education in the United States was accomplished using the mentor system. A physician would finish a rotating internship and move to some town and be taught by the local physicians the skills needed for that particular town. This allowed each community's needs to be met by the teaching of the new general practitioner the skills needed in that community. This also allowed the new physician to start making a living and raising a family, etc. General practitioners would be the surgeons, the obstetricians, and the internists for their given communities. Changes in demographics and the growing complexities of the developing bodies of knowledge made it necessary to produce more highly trained surgeons and other specialists. For many physicians it was a natural desire to want to be considered "specialists".
What was not anticipated by many physicians is that an option to be a generalist would lose its prestige and be further degraded by a growing bureaucracy of insurance and hospitals requiring board certification. It has been shown that there is no statistically significant correlation between board certification and patient safety or quality of care which is why 47 states do not require board certification to practice medicine. Board certification agencies have been increasing their fees exponentially since establishment and the board examinations are known to not be clinically relevant and are at least 5 years out of date. Yet, there is still a misbelief that board certification is necessary to practice medicine and therefore it has made a non-board eligible general physician a rare breed of physician due to the lack of available job opportunities for them.
Certificates of Added Qualifications (CAQs) in adolescent medicine, geriatric medicine, sports medicine, sleep medicine, and hospice and palliative medicine are available for those board-certified family physicians with additional residency training requirements. Recently, new fellowships in International Family Medicine have emerged. These fellowships are designed to train family physicians working in resource-poor environments.
In 2009 an ongoing shortage of primary care physicians (and also other primary care providers) was reported. It was due to several factors, notably the lesser prestige associated with the young specialty, the lower pay, and the increasingly frustrating practice environment. In the US physicians are increasingly forced to do more administrative work, and pay higher malpractice insurance premiums.
Canada
The College of General Practice of Canada was founded in 1954 but in 1967 changed its name to College of Family Physicians of Canada (CFPC).
Oceania
Australia
General Practice in Australia and New Zealand has undergone many changes in training requirements over the past decade. The basic medical degree in Australia is the MBBS (Bachelor of Medicine, Bachelor of Surgery), which has traditionally been attained after completion of an undergraduate five or six-year course. Over the last few years, an ever-increasing number of post-graduate four-year medical programs (previous bachelor's degree required) have become more common and now account more than half of all Australian medical graduates. After graduating, a one-year internship is completed in a public and private hospitals prior to obtaining full registration. Many newly registered medical practitioners undergo one year or more of pre-vocational position as Resident Medical Officers (different titles depending on jurisdictions) before specialist training begins. For general practice training, the medical practitioner then applies to enter a three- or four-year program either through the "Australian General Practice Training Program", "Remote Vocational Training Scheme" or "Independent Pathway". The Australian Government has announced an expansion of the number of GP training places through the AGPT program- 1,500 places per year will be available by 2015.
A combination of coursework and apprenticeship type training leading to the awarding of the FRACGP (Fellowship of the Royal Australian College of General Practitioners) or FACRRM (Fellowship of Australian College of Rural and Remote Medicine), if successful. Since 1996 this qualification or its equivalent has been required in order for new GPs to access Medicare rebates as a specialist general practitioner. Doctors who graduated prior to 1992 and who had worked in general practice for a specified period of time were recognized as "Vocationally Registered" or "VR" GPs, and given automatic and continuing eligibility for general practice Medicare rebates. There is a sizable group of doctors who have identical qualifications and experience, but who have been denied access to VR recognition. They are termed "Non-Vocationally Registered" or so-called "non-VR" GPs. The federal government of Australia recognizes the experience and competence of these doctors, by allowing them access to the "specialist" GP Medicare rebates for working in areas of government policy priority, such as areas of workforce shortage, and metropolitan after hours service. Some programs awarded permanent and unrestricted eligibility for VR rebate levels after 5 years of practice under the program. There is a community-based campaign in support of these so-called Non-VR doctors being granted full and permanent recognition of their experience and expertise, as fully identical with the previous generation of pre-1996 "grandfathered" GPs. This campaign is supported by the official policy of the Australian Medical Association (AMA).
Medicare is Australia's universal health insurance system, and without access to it, a practitioner cannot effectively work in private practice in Australia.
Procedural General Practice training in combination with General Practice Fellowship was first established by the "Australian College of Rural and Remote Medicine" in 2004. This new fellowship was developed in aid to recognise the specialised skills required to work within a rural and remote context. In addition it was hoped to recognise the impending urgency of training Rural Procedural Practitioners to sustain Obstetric and Surgical services within rural Australia. Each training registrar select a speciality that can be used in a rural area from the Advanced Skills Training list and spends a minimum of 12 months completing this specialty, the most common of which are Surgery, Obstetrics/Gynaecology and Anaesthetics. Further choices of specialty include Aboriginal and Torres Strait Islander Health, Adult Internal Medicine, Emergency Medicine, Mental Health, Paediatrics, Population health and Remote Medicine. Shortly after the establishment of the FACRRM, the Royal Australian College of General Practitioners introduced an additional training year (from the basic 3 years) to offer the "Fellowship in Advanced Rural General Practice". The additional year, or Advanced Rural Skills Training (ARST) can be conducted in various locations from Tertiary Hospitals to Small General Practice.
The competent authority pathway is a work-based place assessment process to support International Medical Graduates (IMGs) wishing to work in General Practice. Approval for the ACRRM to undertake these assessments was granted by the Australian Medical Council in August 2010 and the process is to be streamlined in July 2014.
New Zealand
In New Zealand, most GPs work in clinics and health centres usually as part of a Primary Health Organisation (PHO). These are funded at a population level, based on the characteristics of a practice's enrolled population (referred to as capitation-based funding). Fee-for-service arrangements still exist with other funders such as Accident Compensation Corporation (ACC) and Ministry of Social Development (MSD), as well as receiving co-payments from patients to top-up the capitation-based funding.
The basic medical degree in New Zealand is the MBChB degree (Bachelor of Medicine, Bachelor of Surgery), which has traditionally been attained after completion of an undergraduate five or six-year course. In NZ new graduates must complete the GPEP (General Practice Education Program) Stages I and II in order to be granted the title Fellowship of the Royal New Zealand College of General Practitioners (FRNZCGP), which includes the PRIMEX assessment and further CME and Peer group learning sessions as directed by the RNZCGP. Holders of the award of FRNZCGP may apply for specialist recognition with the New Zealand Medical Council (MCNZ), after which they are considered specialists in General Practice by the council and the community. In 2009 the NZ Government increased the number of places available on the state-funded programme for GP training.
There is a shortage of GPs in rural areas and increasingly outer metropolitan areas of large cities, which has led to the use of overseas trained doctors (international medical graduates (IMGs)).
| Biology and health sciences | Health professionals | Health |
197847 | https://en.wikipedia.org/wiki/Imperative%20programming | Imperative programming | In computer science, imperative programming is a programming paradigm of software that uses statements that change a program's state. In much the same way that the imperative mood in natural languages expresses commands, an imperative program consists of commands for the computer to perform. Imperative programming focuses on describing how a program operates step by step (generally order of the steps being determined in source code by the placement of statements one below the other), rather than on high-level descriptions of its expected results.
The term is often used in contrast to declarative programming, which focuses on what the program should accomplish without specifying all the details of how the program should achieve the result.
Procedural programming
Procedural programming is a type of imperative programming in which the program is built from one or more procedures (also termed subroutines or functions). The terms are often used as synonyms, but the use of procedures has a dramatic effect on how imperative programs appear and how they are constructed. Heavy procedural programming, in which state changes are localized to procedures or restricted to explicit arguments and returns from procedures, is a form of structured programming. Since the 1960s, structured programming and modular programming in general have been promoted as techniques to improve the maintainability and overall quality of imperative programs. The concepts behind object-oriented programming attempt to extend this approach.
Procedural programming could be considered a step toward declarative programming. A programmer can often tell, simply by looking at the names, arguments, and return types of procedures (and related comments), what a particular procedure is supposed to do, without necessarily looking at the details of how it achieves its result. At the same time, a complete program is still imperative since it fixes the statements to be executed and their order of execution to a large extent.
Rationale and foundations of imperative programming
The programming paradigm used to build programs for almost all computers typically follows an imperative model. Digital computer hardware is designed to execute machine code, which is native to the computer and is usually written in the imperative style, although low-level compilers and interpreters using other paradigms exist for some architectures such as lisp machines.
From this low-level perspective, the program state is defined by the contents of memory, and the statements are instructions in the native machine language of the computer. Higher-level imperative languages use variables and more complex statements, but still follow the same paradigm. Recipes and process checklists, while not computer programs, are also familiar concepts that are similar in style to imperative programming; each step is an instruction, and the physical world holds the state. Since the basic ideas of imperative programming are both conceptually familiar and directly embodied in the hardware, most computer languages are in the imperative style.
Assignment statements, in imperative paradigm, perform an operation on information located in memory and store the results in memory for later use. High-level imperative languages, in addition, permit the evaluation of complex expressions, which may consist of a combination of arithmetic operations and function evaluations, and the assignment of the resulting value to memory. Looping statements (as in while loops, do while loops, and for loops) allow a sequence of statements to be executed multiple times. Loops can either execute the statements they contain a predefined number of times, or they can execute them repeatedly until some condition is met. Conditional branching statements allow a sequence of statements to be executed only if some condition is met. Otherwise, the statements are skipped and the execution sequence continues from the statement following them. Unconditional branching statements allow an execution sequence to be transferred to another part of a program. These include the jump (called goto in many languages), switch, and the subprogram, subroutine, or procedure call (which usually returns to the next statement after the call).
Early in the development of high-level programming languages, the introduction of the block enabled the construction of programs in which a group of statements and declarations could be treated as if they were one statement. This, alongside the introduction of subroutines, enabled complex structures to be expressed by hierarchical decomposition into simpler procedural structures.
Many imperative programming languages (such as Fortran, BASIC, and C) are abstractions of assembly language.
History of imperative and object-oriented languages
The earliest imperative languages were the machine languages of the original computers. In these languages, instructions were very simple, which made hardware implementation easier but hindered the creation of complex programs. FORTRAN, developed by John Backus at International Business Machines (IBM) starting in 1954, was the first major programming language to remove the obstacles presented by machine code in the creation of complex programs. FORTRAN was a compiled language that allowed named variables, complex expressions, subprograms, and many other features now common in imperative languages. The next two decades saw the development of many other major high-level imperative programming languages. In the late 1950s and 1960s, ALGOL was developed in order to allow mathematical algorithms to be more easily expressed and even served as the operating system's target language for some computers. MUMPS (1966) carried the imperative paradigm to a logical extreme, by not having any statements at all, relying purely on commands, even to the extent of making the IF and ELSE commands independent of each other, connected only by an intrinsic variable named $TEST. COBOL (1960) and BASIC (1964) were both attempts to make programming syntax look more like English. In the 1970s, Pascal was developed by Niklaus Wirth, and C was created by Dennis Ritchie while he was working at Bell Laboratories. Wirth went on to design Modula-2 and Oberon. For the needs of the United States Department of Defense, Jean Ichbiah and a team at Honeywell began designing Ada in 1978, after a 4-year project to define the requirements for the language. The specification was first published in 1983, with revisions in 1995, 2005, and 2012.
The 1980s saw a rapid growth in interest in object-oriented programming. These languages were imperative in style, but added features to support objects. The last two decades of the 20th century saw the development of many such languages. Smalltalk-80, originally conceived by Alan Kay in 1969, was released in 1980, by the Xerox Palo Alto Research Center (PARC). Drawing from concepts in another object-oriented language—Simula (which is considered the world's first object-oriented programming language, developed in the 1960s)—Bjarne Stroustrup designed C++, an object-oriented language based on C. Design of C++ began in 1979 and the first implementation was completed in 1983. In the late 1980s and 1990s, the notable imperative languages drawing on object-oriented concepts were Perl, released by Larry Wall in 1987; Python, released by Guido van Rossum in 1990; Visual Basic and Visual C++ (which included Microsoft Foundation Class Library (MFC) 2.0), released by Microsoft in 1991 and 1993 respectively; PHP, released by Rasmus Lerdorf in 1994; Java, by James Gosling (Sun Microsystems) in 1995, JavaScript, by Brendan Eich (Netscape), and Ruby, by Yukihiro "Matz" Matsumoto, both released in 1995. Microsoft's .NET Framework (2002) is imperative at its core, as are its main target languages, VB.NET and C# that run on it; however Microsoft's F#, a functional language, also runs on it.
Examples
Fortran
FORTRAN (1958) was unveiled as "The IBM Mathematical FORmula TRANslating system." It was designed for scientific calculations, without string handling facilities. Along with declarations, expressions, and statements, it supported:
arrays
subroutines
"do" loops
It succeeded because:
programming and debugging costs were below computer running costs
it was supported by IBM
applications at the time were scientific.
However, non IBM vendors also wrote Fortran compilers, but with a syntax that would likely fail IBM's compiler. The American National Standards Institute (ANSI) developed the first Fortran standard in 1966. In 1978, Fortran 77 became the standard until 1991. Fortran 90 supports:
records
pointers to arrays
COBOL
COBOL (1959) stands for "COmmon Business Oriented Language." Fortran manipulated symbols. It was soon realized that symbols did not need to be numbers, so strings were introduced. The US Department of Defense influenced COBOL's development, with Grace Hopper being a major contributor. The statements were English-like and verbose. The goal was to design a language so managers could read the programs. However, the lack of structured statements hindered this goal.
COBOL's development was tightly controlled, so dialects did not emerge to require ANSI standards. As a consequence, it was not changed for 15 years until 1974. The 1990s version did make consequential changes, like object-oriented programming.
Algol
ALGOL (1960) stands for "ALGOrithmic Language." It had a profound influence on programming language design. Emerging from a committee of European and American programming language experts, it used standard mathematical notation and had a readable structured design. Algol was first to define its syntax using the Backus–Naur form. This led to syntax-directed compilers. It added features like:
block structure, where variables were local to their block
arrays with variable bounds
"for" loops
functions
recursion
Algol's direct descendants include Pascal, Modula-2, Ada, Delphi and Oberon on one branch. On another branch there's C, C++ and Java.
Basic
BASIC (1964) stands for "Beginner's All Purpose Symbolic Instruction Code." It was developed at Dartmouth College for all of their students to learn. If a student did not go on to a more powerful language, the student would still remember Basic. A Basic interpreter was installed in the microcomputers manufactured in the late 1970s. As the microcomputer industry grew, so did the language.
Basic pioneered the interactive session. It offered operating system commands within its environment:
The 'new' command created an empty slate
Statements evaluated immediately
Statements could be programmed by preceding them with a line number
The 'list' command displayed the program
The 'run' command executed the program
However, the Basic syntax was too simple for large programs. Recent dialects added structure and object-oriented extensions. Microsoft's Visual Basic is still widely used and produces a graphical user interface.
C
C programming language (1973) got its name because the language BCPL was replaced with B, and AT&T Bell Labs called the next version "C." Its purpose was to write the UNIX operating system. C is a relatively small language -- making it easy to write compilers. Its growth mirrored the hardware growth in the 1980s. Its growth also was because it has the facilities of assembly language, but uses a high-level syntax. It added advanced features like:
inline assembler
arithmetic on pointers
pointers to functions
bit operations
freely combining complex operators
C allows the programmer to control in which region of memory data is to be stored. Global variables and static variables require the fewest clock cycles to store. The stack is automatically used for the standard variable declarations. Heap memory is returned to a pointer variable from the malloc() function.
The global and static data region is located just above the program region. (The program region is technically called the text region. It's where machine instructions are stored.)
The global and static data region is technically two regions. One region is called the initialized data segment, where variables declared with default values are stored. The other region is called the block started by segment, where variables declared without default values are stored.
Variables stored in the global and static data region have their addresses set at compile-time. They retain their values throughout the life of the process.
The global and static region stores the global variables that are declared on top of (outside) the main() function. Global variables are visible to main() and every other function in the source code.
On the other hand, variable declarations inside of main(), other functions, or within { } block delimiters are local variables. Local variables also include formal parameter variables. Parameter variables are enclosed within the parenthesis of function definitions. They provide an interface to the function.
Local variables declared using the static prefix are also stored in the global and static data region. Unlike global variables, static variables are only visible within the function or block. Static variables always retain their value. An example usage would be the function int increment_counter(){ static int counter = 0; counter++; return counter;}
The stack region is a contiguous block of memory located near the top memory address. Variables placed in the stack are populated from top to bottom. A stack pointer is a special-purpose register that keeps track of the last memory address populated. Variables are placed into the stack via the assembly language PUSH instruction. Therefore, the addresses of these variables are set during runtime. The method for stack variables to lose their scope is via the POP instruction.
Local variables declared without the static prefix, including formal parameter variables, are called automatic variables and are stored in the stack. They are visible inside the function or block and lose their scope upon exiting the function or block.
The heap region is located below the stack. It is populated from the bottom to the top. The operating system manages the heap using a heap pointer and a list of allocated memory blocks. Like the stack, the addresses of heap variables are set during runtime. An out of memory error occurs when the heap pointer and the stack pointer meet.
C provides the malloc() library function to allocate heap memory. Populating the heap with data is an additional copy function. Variables stored in the heap are economically passed to functions using pointers. Without pointers, the entire block of data would have to be passed to the function via the stack.
C++
In the 1970s, software engineers needed language support to break large projects down into modules. One obvious feature was to decompose large projects physically into separate files. A less obvious feature was to decompose large projects logically into abstract datatypes. At the time, languages supported concrete (scalar) datatypes like integer numbers, floating-point numbers, and strings of characters. Concrete datatypes have their representation as part of their name. Abstract datatypes are structures of concrete datatypes — with a new name assigned. For example, a list of integers could be called integer_list.
In object-oriented jargon, abstract datatypes are called classes. However, a class is only a definition; no memory is allocated. When memory is allocated to a class, it's called an object.
Object-oriented imperative languages developed by combining the need for classes and the need for safe functional programming. A function, in an object-oriented language, is assigned to a class. An assigned function is then referred to as a method, member function, or operation. Object-oriented programming is executing operations on objects.
Object-oriented languages support a syntax to model subset/superset relationships. In set theory, an element of a subset inherits all the attributes contained in the superset. For example, a student is a person. Therefore, the set of students is a subset of the set of persons. As a result, students inherit all the attributes common to all persons. Additionally, students have unique attributes that other persons don't have. Object-oriented languages model subset/superset relationships using inheritance. Object-oriented programming became the dominant language paradigm by the late 1990s.
C++ (1985) was originally called "C with Classes." It was designed to expand C's capabilities by adding the object-oriented facilities of the language Simula.
An object-oriented module is composed of two files. The definitions file is called the header file. Here is a C++ header file for the GRADE class in a simple school application:
// grade.h
// -------
// Used to allow multiple source files to include
// this header file without duplication errors.
// See: https://en.wikipedia.org/wiki/Include_guard
// ----------------------------------------------
#ifndef GRADE_H
#define GRADE_H
class GRADE {
public:
// This is the constructor operation.
// ----------------------------------
GRADE ( const char letter );
// This is a class variable.
// -------------------------
char letter;
// This is a member operation.
// ---------------------------
int grade_numeric( const char letter );
// This is a class variable.
// -------------------------
int numeric;
};
#endif
A constructor operation is a function with the same name as the class name. It is executed when the calling operation executes the new statement.
A module's other file is the source file. Here is a C++ source file for the GRADE class in a simple school application:
// grade.cpp
// ---------
#include "grade.h"
GRADE::GRADE( const char letter )
{
// Reference the object using the keyword 'this'.
// ----------------------------------------------
this->letter = letter;
// This is Temporal Cohesion
// -------------------------
this->numeric = grade_numeric( letter );
}
int GRADE::grade_numeric( const char letter )
{
if ( ( letter == 'A' || letter == 'a' ) )
return 4;
else
if ( ( letter == 'B' || letter == 'b' ) )
return 3;
else
if ( ( letter == 'C' || letter == 'c' ) )
return 2;
else
if ( ( letter == 'D' || letter == 'd' ) )
return 1;
else
if ( ( letter == 'F' || letter == 'f' ) )
return 0;
else
return -1;
}
Here is a C++ header file for the PERSON class in a simple school application:
// person.h
// --------
#ifndef PERSON_H
#define PERSON_H
class PERSON {
public:
PERSON ( const char *name );
const char *name;
};
#endif
Here is a C++ source file for the PERSON class in a simple school application:
// person.cpp
// ----------
#include "person.h"
PERSON::PERSON ( const char *name )
{
this->name = name;
}
Here is a C++ header file for the STUDENT class in a simple school application:
// student.h
// ---------
#ifndef STUDENT_H
#define STUDENT_H
#include "person.h"
#include "grade.h"
// A STUDENT is a subset of PERSON.
// --------------------------------
class STUDENT : public PERSON{
public:
STUDENT ( const char *name );
~STUDENT();
GRADE *grade;
};
#endif
Here is a C++ source file for the STUDENT class in a simple school application:
// student.cpp
// -----------
#include "student.h"
#include "person.h"
STUDENT::STUDENT ( const char *name ):
// Execute the constructor of the PERSON superclass.
// -------------------------------------------------
PERSON( name )
{
// Nothing else to do.
// -------------------
}
STUDENT::~STUDENT()
{
// deallocate grade's memory
// to avoid memory leaks.
// -------------------------------------------------
delete this->grade;
}
Here is a driver program for demonstration:
// student_dvr.cpp
// ---------------
#include <iostream>
#include "student.h"
int main( void )
{
STUDENT *student = new STUDENT( "The Student" );
student->grade = new GRADE( 'a' );
std::cout
// Notice student inherits PERSON's name
<< student->name
<< ": Numeric grade = "
<< student->grade->numeric
<< "\n";
// deallocate student's memory
// to avoid memory leaks.
// -------------------------------------------------
delete student;
return 0;
}
Here is a makefile to compile everything:
# makefile
# --------
all: student_dvr
clean:
rm student_dvr *.o
student_dvr: student_dvr.cpp grade.o student.o person.o
c++ student_dvr.cpp grade.o student.o person.o -o student_dvr
grade.o: grade.cpp grade.h
c++ -c grade.cpp
student.o: student.cpp student.h
c++ -c student.cpp
person.o: person.cpp person.h
c++ -c person.cpp
| Technology | Software development: General | null |
197964 | https://en.wikipedia.org/wiki/Electron%20affinity | Electron affinity | The electron affinity (Eea) of an atom or molecule is defined as the amount of energy released when an electron attaches to a neutral atom or molecule in the gaseous state to form an anion.
X(g) + e− → X−(g) + energy
This differs by sign from the energy change of electron capture ionization. The electron affinity is positive when energy is released on electron capture.
In solid state physics, the electron affinity for a surface is defined somewhat differently (see below).
Measurement and use of electron affinity
This property is used to measure atoms and molecules in the gaseous state only, since in a solid or liquid state their energy levels would be changed by contact with other atoms or molecules.
A list of the electron affinities was used by Robert S. Mulliken to develop an electronegativity scale for atoms, equal to the average of the electrons
affinity and ionization potential. Other theoretical concepts that use electron affinity include electronic chemical potential and chemical hardness. Another example, a molecule or atom that has a more positive value of electron affinity than another is often called an electron acceptor and the less positive an electron donor. Together they may undergo charge-transfer reactions.
Sign convention
To use electron affinities properly, it is essential to keep track of sign. For any reaction that releases energy, the change ΔE in total energy has a negative value and the reaction is called an exothermic process. Electron capture for almost all non-noble gas atoms involves the release of energy and thus is exothermic. The positive values that are listed in tables of Eea are amounts or magnitudes. It is the word "released" within the definition "energy released" that supplies the negative sign to ΔE. Confusion arises in mistaking Eea for a change in energy, ΔE, in which case the positive values listed in tables would be for an endo- not exo-thermic process. The relation between the two is Eea = −ΔE(attach).
However, if the value assigned to Eea is negative, the negative sign implies a reversal of direction, and energy is required to attach an electron. In this case, the electron capture is an endothermic process and the relationship, Eea = −ΔE(attach) is still valid. Negative values typically arise for the capture of a second electron, but also for the nitrogen atom.
The usual expression for calculating Eea when an electron is attached is
This expression does follow the convention ΔX = X(final) − X(initial) since −ΔE = −(E(final) − E(initial)) = E(initial) − E(final).
Equivalently, electron affinity can also be defined as the amount of energy required to detach an electron from the atom while it holds a single-excess-electron thus making the atom a negative ion, i.e. the energy change for the process
X− → X + e−
If the same table is employed for the forward and reverse reactions, without switching signs, care must be taken to apply the correct definition to the corresponding direction, attachment (release) or detachment (require). Since almost all detachments (require +) an amount of energy listed on the table, those detachment reactions are endothermic, or ΔE(detach) > 0.
Electron affinities of the elements
Although Eea varies greatly across the periodic table, some patterns emerge. Generally, nonmetals have more positive Eea than metals. Atoms whose anions are more stable than neutral atoms have a greater Eea. Chlorine most strongly attracts extra electrons; neon most weakly attracts an extra electron. The electron affinities of the noble gases have not been conclusively measured, so they may or may not have slightly negative values.
Eea generally increases across a period (row) in the periodic table prior to reaching group 18. This is caused by the filling of the valence shell of the atom; a group 17 atom releases more energy than a group 1 atom on gaining an electron because it obtains a filled valence shell and therefore is more stable. In group 18, the valence shell is full, meaning that added electrons are unstable, tending to be ejected very quickly.
Counterintuitively, Eea does not decrease when progressing down most columns of the periodic table. For example, Eea actually increases consistently on descending the column for the group 2 data. Thus, electron affinity follows the same "left-right" trend as electronegativity, but not the "up-down" trend.
The following data are quoted in kJ/mol.
Molecular electron affinities
The electron affinity of molecules is a complicated function of their electronic structure.
For instance the electron affinity for benzene is negative, as is that of naphthalene, while those of anthracene, phenanthrene and pyrene are positive. In silico experiments show that the electron affinity of hexacyanobenzene surpasses that of fullerene.
"Electron affinity" as defined in solid state physics
In the field of solid state physics, the electron affinity is defined differently than in chemistry and atomic physics. For a semiconductor-vacuum interface (that is, the surface of a semiconductor), electron affinity, typically denoted by EEA or χ, is defined as the energy obtained by moving an electron from the vacuum just outside the semiconductor to the bottom of the conduction band just inside the semiconductor:
In an intrinsic semiconductor at absolute zero, this concept is functionally analogous to the chemistry definition of electron affinity, since an added electron will spontaneously go to the bottom of the conduction band. At nonzero temperature, and for other materials (metals, semimetals, heavily doped semiconductors), the analogy does not hold since an added electron will instead go to the Fermi level on average. In any case, the value of the electron affinity of a solid substance is very different from the chemistry and atomic physics electron affinity value for an atom of the same substance in gas phase. For example, a silicon crystal surface has electron affinity 4.05 eV, whereas an isolated silicon atom has electron affinity 1.39 eV.
The electron affinity of a surface is closely related to, but distinct from, its work function. The work function is the thermodynamic work that can be obtained by reversibly and isothermally removing an electron from the material to vacuum; this thermodynamic electron goes to the Fermi level on average, not the conduction band edge: . While the work function of a semiconductor can be changed by doping, the electron affinity ideally does not change with doping and so it is closer to being a material constant. However, like work function the electron affinity does depend on the surface termination (crystal face, surface chemistry, etc.) and is strictly a surface property.
In semiconductor physics, the primary use of the electron affinity is not actually in the analysis of semiconductor–vacuum surfaces, but rather in heuristic electron affinity rules for estimating the band bending that occurs at the interface of two materials, in particular metal–semiconductor junctions and semiconductor heterojunctions.
In certain circumstances, the electron affinity may become negative. Often negative electron affinity is desired to obtain efficient cathodes that can supply electrons to the vacuum with little energy loss. The observed electron yield as a function of various parameters such as bias voltage or illumination conditions can be used to describe these structures with band diagrams in which the electron affinity is one parameter. For one illustration of the apparent effect of surface termination on electron emission, see Figure 3 in Marchywka Effect.
| Physical sciences | Periodic table | Chemistry |
198030 | https://en.wikipedia.org/wiki/Pont%20du%20Gard | Pont du Gard | The Pont du Gard is an ancient Roman aqueduct bridge built in the first century AD to carry water over to the Roman colony of Nemausus (Nîmes). It crosses the river Gardon near the town of Vers-Pont-du-Gard in southern France. The Pont du Gard is one of the best preserved Roman aqueduct bridges. It was added to UNESCO's list of World Heritage sites in 1985 because of its exceptional preservation, historical importance, and architectural ingenuity.
Description
The bridge has three tiers of arches made from Shelly limestone and stands high. The aqueduct formerly carried an estimated of water a day over to the fountains, baths and homes of the citizens of Nîmes. The structure's precise construction allowed an average gradient of in . It may have been in use as late as the 6th century, with some parts used for significantly longer, but lack of maintenance after the 4th century led to clogging by mineral deposits and debris that eventually stopped the flow of water.
After the Roman Empire collapsed and the aqueduct fell into disuse, the Pont du Gard remained largely intact with a secondary function as a toll bridge. For centuries the local lords and bishops were responsible for its upkeep, with a right to levy tolls on travellers using it to cross the river. Over time, some of its stone blocks were looted, and serious damage was inflicted in the 17th century. It attracted increasing attention starting in the 18th century, and became an important tourist destination. A series of renovations between the 18th and 21st centuries, commissioned by local authorities and the French state, culminated in 2000 with opening of a new visitor centre and removal of traffic and buildings from the bridge and area immediately around it. Today it is one of France's most popular tourist attractions, and has attracted the attention of a succession of literary and artistic visitors.
Route of the Nîmes aqueduct
The location of Nemausus (Nîmes) was somewhat inconvenient when it came to providing a water supply. Plains lie to the city's south and east, where any sources of water would be at too low an altitude to be able to flow to the city, while the hills to the west made a water supply route too difficult from an engineering point of view. The only real alternative was to look to the north and in particular to the area around Ucetia (Uzès), where there are natural springs.
The Nîmes aqueduct was built to channel water from the springs of the Fontaine d'Eure near Uzès to the castellum divisorum (repartition basin) in Nemausus. From there, it was distributed to fountains, baths and private homes around the city. The straight-line distance between the two is only about , but the aqueduct takes a winding route measuring around . This was necessary to circumvent the southernmost foothills of the Massif Central, known as the . They are difficult to cross, as they are covered in dense vegetation and garrigue and indented by deep valleys. It was impractical for the Romans to attempt to tunnel through the hills, as it would have required a tunnel of between , depending on the starting point. A roughly V-shaped course around the eastern end of the Garrigues de Nîmes was therefore the only practical way of transporting the water from the spring to the city.
The Fontaine d'Eure, at above sea level, is only higher than the repartition basin in Nîmes, but this provided a sufficient gradient to sustain a steady flow of water to the 50,000 inhabitants of the Roman city. The aqueduct's average gradient is only 1 in 3,000. It varies widely along its course, but is as little as 1 in 20,000 in some sections. The Pont du Gard itself descends in , a gradient of 1 in 18,241. The average gradient between the start and end of the aqueduct is far shallower than was usual for Roman aqueducts – only about a tenth of the average gradient of some of the aqueducts in Rome.
The reason for the disparity in gradients along the aqueduct's route is that a uniform gradient would have meant that the Pont du Gard would have been infeasibly high, given the limitations of the technology of the time. By varying the gradient along the route, the aqueduct's engineers were able to lower the height of the bridge by to above the river – still exceptionally high by Roman standards, but within acceptable limits. This height limit governed the profile and gradients of the entire aqueduct, but it came at the price of creating a "sag" in the middle of the aqueduct. The gradient profile before the Pont du Gard is relatively steep, descending at per kilometre, but thereafter it descends by only over the remaining . In one section, the winding route between the Pont du Gard and St Bonnet required an extraordinary degree of accuracy from the Roman engineers, who had to allow for a fall of only per of the conduit.
It is estimated that the aqueduct supplied the city with around of water a day that took nearly 27 hours to flow from the source to the city. The water arrived in the castellum divisorum at Nîmes – an open, shallow, circular basin 5.5 m in diameter by 1 m deep. It would have been surrounded by a balustrade within some sort of enclosure, probably under some kind of small but elaborate pavilion. When it was excavated, traces of a tiled roof, Corinthian columns and a fresco decorated with fish and dolphins were discovered in a fragmentary condition. The aqueduct water entered through an opening wide, and ten large holes in the facing wall, each wide, directed the water into the city's main water pipes. Three large drains were also located in the floor, possibly to enable the nearby amphitheatre to be flooded rapidly to enable naumachia (mock naval battles) to be held.
The spring still exists and is now the site of a small modern pumping station. Its water is pure but high in dissolved calcium carbonate leached out of the surrounding limestone. This presented the Romans with significant problems in maintaining the aqueduct, as the carbonates precipitated out of the water during its journey through the conduit. This caused the flow of the aqueduct to become progressively reduced by deposits of calcareous sinter. Another threat was posed by vegetation penetrating the stone lid of the channel. As well as obstructing the flow of the water, dangling roots introduced algae and bacteria that decomposed in a process called biolithogenesis, producing concretions within the conduit. It required constant maintenance by circitores, workers responsible for the aqueduct's upkeep, who crawled along the conduit scrubbing the walls clean and removing any vegetation.
Much of the Nîmes aqueduct was built underground, as was typical of Roman aqueducts. It was constructed by digging a trench in which a stone channel was built and enclosed by an arched roof of stone slabs, which was then covered with earth. Some sections of the channel are tunnelled through solid rock. In all, of the aqueduct was constructed below the ground. The remainder had to be carried on the surface through conduits set on a wall or on arched bridges. Some substantial remains of the above-ground works can still be seen today, such as the so-called "Pont Rue" that stretches for hundreds of metres around Vers and still stands up to high. Other surviving parts include the Pont de Bornègre, three arches carrying the aqueduct across a stream; the Pont de Sartanette, near the Pont du Gard, which covers across a small valley; and three sections of aqueduct tunnel near Sernhac, measuring up to long. However, the Pont du Gard is by far the best-preserved section of the entire aqueduct.
Description of the bridge
Built on three levels, the Pont is high above the river at low water and long. Its width varies from at the bottom to at the top. The three levels of arches are recessed, with the main piers in line one above another. The span of the arches varies slightly, as each was constructed independently to provide flexibility to protect against subsidence. Each level has a differing number of arches:
The first level of the Pont du Gard adjoins a road bridge that was added in the 18th century. The water conduit or specus, which is about high and wide, is carried at the top of the third level. The upper levels of the bridge are slightly curved in the upstream direction. It was long believed that the engineers had designed it this way deliberately to strengthen the bridge's structure against the flow of water, like a dam wall. However, a microtopographic survey carried out in 1989 showed that the bend is caused by the stone expanding and contracting by about a day under the heat of the sun. Over the centuries, this process has produced the current deformation.
The Pont du Gard was constructed largely without the use of mortar or clamps. It contains an estimated 50,400 tons of limestone with a volume of some ; some of the individual blocks weigh up to 6 tons. Most of the stone was extracted from the local quarry of Estel located approximately downstream, on the banks of the Gardon River. The coarse-grained soft reddish shelly limestone, known locally as "Pierre de Vers", lends itself very well to dimension stone production. The blocks were precisely cut to fit perfectly together by friction and gravity, eliminating the need for mortar. The builders also left inscriptions on the stonework conveying various messages and instructions. Many blocks were numbered and inscribed with the required locations, such as fronte dextra or fronte sinistra (front right or front left), to guide the builders.
The method of construction is fairly well understood by historians. The patron of the aqueduct – a rich individual or the city of Nîmes itself – would have hired a large team of contractors and skilled labourers. A surveyor or mensor planned the route using a groma for sighting, the chorobates for levelling, and a set of measuring poles five or ten Roman feet long. His figures and perhaps diagrams were recorded on wax tablets, later to be written up on scrolls. The builders may have used templates to guide them with tasks that required a high degree of precision, such as carving the standardised blocks from which the water conduit was constructed.
The builders would have made extensive use of cranes and block and tackle pulleys to lift the stones into place. Much of the work could have been done using simple sheers operated by a windlass. For the largest blocks, a massive human-powered treadmill would have been used; such machines were still being used in the quarries of Provence until as late as the start of the 20th century. A complex scaffold was erected to support the bridge as it was being built. Large blocks were left protruding from the bridge to support the frames and scaffolds used during construction. The aqueduct as a whole would have been a very expensive undertaking; Émile Espérandieu estimated the cost to be over 30 million sesterces, equivalent to 50 years' pay for 500 new recruits in a Roman legion.
Although the exterior of the Pont du Gard is rough and relatively unfinished, the builders took care to ensure that the interior of the water conduit was as smooth as possible so that the flow of water would not be obstructed. The walls of the conduit were constructed from dressed masonry and the floor from concrete. Both were covered with a stucco incorporating minute shards of pottery and tile. It was painted with olive oil and covered with maltha, a mixture of slaked lime, pork grease and the viscous juice of unripe figs. This produced a surface that was both smooth and durable.
Although the Pont du Gard is renowned for its appearance, its design is not optimal as the technique of stacking arches on top of each other is clumsy and inefficient (and therefore expensive) in the amount of materials it requires. Later aqueducts had a more sophisticated design, making greater use of concrete to reduce their volume and cost of construction. The Aqueduct bridge of Segovia and the Pont de les Ferreres are of roughly similar length but use far fewer arches. Roman architects were eventually able to do away with "stacking" altogether. The Acueducto de los Milagros in Mérida, Spain and the Chabet Ilelouine aqueduct bridge, near Cherchell, Algeria utilise tall, slender piers, constructed from top to bottom with concrete-faced masonry and brick.
History
The construction of the aqueduct has long been credited to the Roman emperor Augustus' son-in-law and aide, Marcus Vipsanius Agrippa, around the year 19 BC. At the time, he was serving as aedile, the senior magistrate responsible for managing the water supply of Rome and its colonies. Espérandieu, writing in 1926, linked the construction of the aqueduct with Agrippa's visit to Narbonensis in that year. Newer excavations suggest the construction may have taken place between 40 and 60 AD. Tunnels dating from the time of Augustus had to be bypassed by the builders of the Nîmes aqueduct, and coins discovered in the outflow in Nîmes are no older than the reign of the emperor Claudius (41–54 AD). On this basis, a team led by Guilhem Fabre has argued that the aqueduct must have been completed around the middle of the 1st century AD. It is believed to have taken about fifteen years to build, employing between 800 and 1,000 workers.
From the 4th century onwards, the aqueduct's maintenance was neglected as successive waves of invaders disrupted the region. It became clogged with debris, encrustations and plant roots, greatly reducing the flow of the water. The resulting deposits in the conduit, consisting of layers of dirt and organic material, are up to thick on each wall. An analysis of the deposits originally suggested that it had continued to supply water to Nîmes until as late as the 9th century, but more recent investigations suggest that it had gone out of use by about the sixth century, though parts of it may have continued to be used for significantly longer.
Although some of its stones were plundered for use elsewhere, the Pont du Gard remained largely intact. Its survival was due to its use as a toll bridge across the valley. In the 13th century the French king granted the seigneurs of Uzès the right to levy tolls on those using the bridge. The right later passed to the Bishops of Uzès. In return, they were responsible for maintaining the bridge in good repair. However, it suffered serious damage during the 1620s when Henri, Duke of Rohan made use of the bridge to transport his artillery during the wars between the French royalists and the Huguenots, whom he led. To make space for his artillery to cross the bridge, the duke had one side of the second row of arches cut away to a depth of about one-third of their original thickness. This left a gap on the lowest deck wide enough to accommodate carts and cannons, but severely weakened the bridge in the process.
In 1703 the local authorities renovated the Pont du Gard to repair cracks, fill in ruts and replace the stones lost in the previous century. A new bridge was built by the engineer Henri Pitot in 1743–47 next to the arches of the lower level, so that the road traffic could cross on a purpose-built bridge. The novelist Alexandre Dumas was strongly critical of the construction of the new bridge, commenting that "it was reserved for the eighteenth century to dishonour a monument which the barbarians of the fifth had not dared to destroy." The Pont du Gard continued to deteriorate and by the time Prosper Mérimée saw it in 1835 it was at serious risk of collapse from erosion and the loss of stonework.
Napoleon III, who had a great admiration for all things Roman, visited the Pont du Gard in 1850 and took a close interest in it. He approved plans by the architect Charles Laisné to repair the bridge in a project which was carried out between 1855 and 1858, with funding provided by the Ministry of State. The work involved substantial renovations that included replacing the eroded stone, infilling some of the piers with concrete to aid stability and improving drainage by separating the bridge from the aqueduct. Stairs were installed at one end and the conduit walls were repaired, allowing visitors to walk along the conduit itself in reasonable safety.
There have been a number of subsequent projects to consolidate the piers and arches of the Pont du Gard. It has survived three serious floods over the last century; in 1958 the whole of the lower tier was submerged by a giant flood that washed away other bridges, and in 1998 another major flood affected the area. A further flood struck in 2002, badly damaging nearby installations.
The Pont du Gard was added to UNESCO's list of World Heritage Sites in 1985 on the criteria of "Human creative genius; testimony to cultural tradition; significance to human history". The description on the list states: "The hydraulic engineers and ... architects who conceived this bridge created a technical as well as artistic masterpiece."
Tourism
The Pont du Gard has been a tourist attraction for centuries. The outstanding quality of the bridge's masonry led to it becoming an obligatory stop for French journeyman masons on their traditional tour around the country (see Compagnons du Tour de France), many of whom have left their names on the stonework. From the 18th century onwards, particularly after the construction of the new road bridge, it became a famous staging-post for travellers on the Grand Tour and became increasingly renowned as an object of historical importance and French national pride.
The bridge has had a long association with French monarchs seeking to associate themselves with a symbol of Roman imperial power. King Charles IX of France visited in 1564 during his Grand Tour of France and was greeted with a grand entertainment laid on by the Duc d'Uzès. Twelve young girls dressed as nymphs came out of a cave by the riverside near the aqueduct and presented the king with pastry and preserved fruits. A century later, Louis XIV and his court visited the Pont du Gard during a visit to Nîmes in January 1660 shortly after the signature of the Treaty of the Pyrenees. In 1786 his great-great-great-grandson Louis XVI commissioned the artist Hubert Robert to produce a set of paintings of Roman ruins of southern France to hang in the king's new dining room at the Palace of Fontainebleau, including a picture depicting the Pont du Gard in an idealised landscape. The commission was meant to reassert the ties between the French monarchy and the imperial past. Napoleon III, in the mid-19th century, consciously identified with Augustus and accorded great respect to Roman antiquities; his patronage of the bridge's restoration in the 1850s was essential to its survival.
By the 1990s the Pont du Gard had become a hugely popular tourist attraction but was congested with traffic – vehicles were still allowed to drive over the 1743 road bridge – and was cluttered with illegally built structures and tourist shops lining the river banks. As the architect Jean-Paul Viguier put it, the "appetite for gain" had transformed the Pont du Gard into "a fairground attraction". In 1996 the General Council of the Gard département began a major four-year project to improve the area, sponsored by the French government, in conjunction with local sources, UNESCO and the EU. The entire area around the bridge was pedestrianised and a new visitor centre was built on the north bank to a design by Jean-Paul Viguier. The redevelopment has ensured that the area around the Pont du Gard is now much quieter due to the removal of vehicle traffic, and the new museum provides a much improved historical context for visitors. The Pont du Gard is today one of France's top five tourist attractions, with 1.4 million visitors reported in 2001.
Literary visitors
Since it became a tourist destination, many novelists and writers have visited the Pont du Gard and written of the experience. Jean-Jacques Rousseau was overwhelmed when he first visited it in 1738:
The novelist Henry James, visiting in 1884, was similarly impressed; he described the Pont du Gard as "unspeakably imposing, and nothing could well be more Roman." He commented:
The mid-19th-century writer Joseph Méry wrote in his 1853 book Les Nuits italiennes, contes nocturnes that on seeing the Pont du Gard:
Hilaire Belloc wrote in 1928 that:
| Technology | Bridges | null |
198041 | https://en.wikipedia.org/wiki/Crane%20%28bird%29 | Crane (bird) | Cranes are a type of large bird with long legs and necks in the biological family Gruidae of the order Gruiformes. The family has 15 species placed in four genera which are Antigone, Balearica, Leucogeranus, and Grus. They are large birds with long necks and legs, a tapering form, and long secondary feathers on the wing that project over the tail. Most species have muted gray or white plumages, marked with black, and red bare patches on the face, but the crowned cranes of the genus Balearica have vibrantly-coloured wings and golden "crowns" of feathers. Cranes fly with their necks extended outwards instead of bent into an S-shape and their long legs outstretched.
Cranes live on most continents, with the exception of Antarctica and South America. Some species and populations of cranes migrate over long distances; others do not migrate at all. Cranes are solitary during the breeding season, occurring in pairs, but during the non-breeding season, most species are gregarious, forming large flocks where their numbers are sufficient.
They are opportunistic feeders that change their diets according to the season and their own nutrient requirements. They eat a range of items from small rodents, eggs of birds, fish, amphibians, and insects to grain and berries. Cranes construct platform nests in shallow water, and typically lay a clutch of two eggs at a time. Both parents help to rear the young, which remain with them until the next breeding season. Most species of cranes have been affected by human activities and are at the least classified as threatened, if not critically endangered. The plight of the whooping cranes of North America inspired some of the first US legislation to protect endangered species.
Description
Cranes are very large birds, often considered the world's tallest flying birds. They range in size from the demoiselle crane, which measures in length, to the sarus crane, which can be up to , although the heaviest is the red-crowned crane, which can weigh prior to migrating. They are long-legged and long-necked birds with streamlined bodies and large, rounded wings. The males and females do not vary in external appearance, but males tend to be slightly larger than females.
The plumage of cranes varies by habitat. Species inhabiting vast, open wetlands tend to have more white in their plumage than do species that inhabit smaller wetlands or forested habitats, which tend to be more grey. These white species are also generally larger. The smaller size and colour of the forest species is thought to help them maintain a less conspicuous profile while nesting; two of these species (the common and sandhill cranes) also daub their feathers with mud which some observers suspect helps them to hide while nesting.
Most crane species have bare patches of skin on their heads and can expand the patches in order to communicate aggression. Species lacking these bare patches use specialized feather tufts to signal similar information.
Also important to communication is the position and length of the trachea. In the two crowned cranes, the trachea is shorter and only slightly impressed upon the bone of the sternum, whereas the trachea of the other species are longer and penetrate the sternum. In some species, the entire sternum is fused to the bony plates of the trachea, and this helps amplify the crane's calls, allowing them to carry for several kilometres.
Taxonomy and systematics
The family name Gruidae comes from the genus Grus, this genus name is obtained from the epithet of the common crane which is Ardea grus, it is named by Carl Linnaeus from the Latin word grus meaning "crane".
The 15 living species of cranes are placed in four genera. A molecular phylogenetic study published in 2010 found that the genus Grus, as then defined, was polyphyletic. In the resulting rearrangement to create monophyletic genera, the Siberian crane was moved to the resurrected monotypic genus Leucogeranus, while the sandhill crane, the white-naped crane, the sarus crane, and the brolga were moved to the resurrected genus Antigone. Some authorities recognize the additional genera Anthropoides (for the demoiselle crane and blue crane) and Bugeranus (for the wattled crane).
The following cladogram is based on a molecular phylogenetic study by Carey Krajewski and collaborators that was published in 2010.
Evolution
The fossil record of cranes is incomplete. Apparently, the subfamilies were well distinct by the Late Eocene (around 35 mya). The present genera are apparently some 20 mya old. Biogeography of known fossil and the living taxa of cranes suggests that the group is probably of (Laurasian?) Old World origin. The extant diversity at the genus level is centered on (eastern) Africa, although no fossil record exists from there. On the other hand, it is peculiar that numerous fossils of Ciconiiformes are documented from there; these birds presumably shared much of their habitat with cranes back then already. Cranes are sister taxa to Eogruidae, a lineage of flightless birds; as predicted by the fossil record of true cranes, eogruids were native to the Old World. A species of true crane, Antigone cubensis, has similarly become flightless and ratite-like.
Fossil genera are tentatively assigned to the present-day subfamilies:
Gruinae
Palaeogrus (Middle Eocene of Germany and Italy – Middle Miocene of France)
Pliogrus (Early Pliocene of Eppelsheim, Germany)
Camusia (Late Miocene of Menorca, Mediterranean)
"Grus" conferta (Late Miocene/Early Pliocene of Contra Costa County, US)
Sometimes considered Balearicinae
Geranopsis (Hordwell Late Eocene – Early Oligocene of England)
Anserpica (Late Oligocene of France)
Sometimes considered Gruidae incertae sedis
Eobalearica (Ferghana Late? Eocene of Ferghana, Uzbekistan)
Probalearica (Late Oligocene? – Middle Pliocene of Florida, US, France?, Moldavia and Mongolia) – A nomen dubium?
Aramornis (Sheep Creek Middle Miocene of Snake Creek Quarries, US)
Distribution and habitat
The cranes have a cosmopolitan distribution, occurring across most of the world continents. They are absent from Antarctica and, mysteriously, South America. East Asia has the highest crane diversity, with eight species, followed by Africa, which is home to five resident species and wintering populations of a sixth. Australia, Europe, and North America have two regularly occurring species each. Of the four crane genera, Balearica (two species) is restricted to Africa, and Leucogeranus (one species) is restricted to Asia; the other two genera, Grus (including Anthropoides and Bugeranus) and Antigone, are both widespread.
Many species of cranes are dependent on wetlands and grasslands, and most species nest in shallow wetlands. Some species nest in wetlands, but move their chicks up onto grasslands or uplands to feed (while returning to wetlands at night), whereas others remain in wetlands for the entirety of the breeding season. Even the demoiselle crane and blue crane, which may nest and feed in grasslands (or even arid grasslands or deserts), require wetlands for roosting at night. The Sarus Crane in south Asia is unique in having a significant breeding population using agricultural fields to breed in areas alongside very high density of humans and intensive farming, largely due to the positive attitudes of farmers towards the cranes. In Australia, the Brolga occurs in the breeding areas of Sarus Cranes in Queensland state, and they achieve sympatry by using different habitats. Sarus Cranes in Queensland largely live in Eucalyptus-dominated riverine, while most Brolgas use non-wooded regional ecosystems that include vast grassland habitats. The only two species that do not always roost in wetlands are the two African crowned cranes (Balearica), which are the only cranes to roost in trees.
Some crane species are sedentary, remaining in the same area throughout the year, while others are highly migratory, traveling thousands of kilometres each year from their breeding sites. A few species like Sarus Cranes have both migratory and sedentary populations, and healthy sedentary populations have a large proportion of cranes that are not territorial, breeding pairs.
Behaviour and ecology
The cranes are diurnal birds that vary in their sociality by season and location. During the breeding season, they are territorial and usually remain on their territory all the time. In contrast in the non-breeding season, they tend to be gregarious, forming large flocks to roost, socialize, and in some species feed. Sarus Crane breeding pairs maintain territories throughout the year in south Asia, and non-breeding birds live in flocks that can also be seen throughout the year. Large aggregations of cranes likely increase safety for individual cranes when resting and flying and also increase chances for young unmated birds to meet partners.
Calls and communication
Cranes are highly vocal and have several specialized calls. The vocabulary begins soon after hatching with low, purring calls for maintaining contact with their parents, as well as food-begging calls. Other calls used as chicks include alarm calls and "flight intention" calls, both of which are maintained into adulthood. Cranes are noticed the most due to their loud duet calls that can be used to distinguish individual pairs. Sarus crane trios produce synchronized unison calls called "triets" whose structure is identical to duets of normal pairs, but have a lower frequency.
Feeding
The cranes consume a wide range of food, both animal and plant matter. When feeding on land, they consume seeds, leaves, nuts and acorns, berries, fruit, insects, worms, snails, small reptiles, mammals, and birds. In wetlands and agriculture fields, roots, rhizomes, tubers, and other parts of emergent plants, other molluscs, small fish, eggs of birds and amphibians are also consumed, as well. The exact composition of the diet varies by location, season, and availability. Within the wide range of items consumed, some patterns are suggested but require specific investigation to confirm; the shorter-billed species usually feed in drier uplands, while the longer-billed species feed in wetlands.
Cranes employ different foraging techniques for different food types and in different habitats. Tubers and rhizomes are dug for and a crane digging for them remains in place for some time digging and then expanding a hole to prise them out of the soil. In contrast both to this and the stationary wait and watch hunting methods employed by many herons, they forage for insects and animal prey by slowly moving forwards with their heads lowered and probing with their bills.
Where more than one species of cranes exists in a locality, each species adopts separate niches to minimise competition. At one important lake in Jiangxi Province in China, the Siberian cranes feed on the mudflats and in shallow water, the white-naped cranes on the wetland borders, the hooded cranes on sedge meadows, and the last two species also feed on the agricultural fields along with the common cranes. In Australia, where Sarus Cranes live alongside Brolgas, they have different diets: Sarus Cranes' diet consisted of diverse vegetation, while Brolga diet spanned a much wider range of trophic levels. Some crane species such as the Common/ Eurasian crane use a kleptoparasitic strategy to recover from temporary reductions in feeding rate, particularly when the rate is below the threshold of intake necessary for survival. Accumulated intake of during daytime shows a typical anti-sigmoid shape, with greatest increases of intake after dawn and before dusk.
Breeding
Cranes are perennially monogamous breeders, establishing long-term pair bonds that may last the lifetime of the birds. Pair bonds begin to form in the second or third years of life, but several years pass before the first successful breeding season. Initial breeding attempts often fail, and in many cases, newer pair bonds dissolve (divorce) after unsuccessful breeding attempts. Pairs that are repeatedly successful at breeding remain together for as long as they continue to do so. In a study of sandhill cranes in Florida, seven of the 22 pairs studied remained together for an 11-year period. Of the pairs that separated, 53% was due to the death of one of the pair, 18% was due to divorce, and the fate of 29% of pairs was unknown. Similar results had been found by acoustic monitoring (sonography/frequency analysis of duet and guard calls) in three breeding areas of common cranes in Germany over 10 years.
Cranes are territorial and generally seasonal breeders. Seasonality varies both between and within species, depending on local conditions. Migratory species begin breeding upon reaching their summer breeding grounds, between April and June. The breeding season of tropical species is usually timed to coincide with the wet or monsoon seasons. Artificial sources of water such as irrigation canals and irregular rainfall can sometimes provide adequate moisture to maintain wetland habitat outside the normal wet season, and allows for occasional aseasonal nesting throughout the year in few tropical species.
Territory sizes also vary depending on location. Tropical species can maintain very small territories, for example sarus cranes in India can breed on territories as small as one hectare where the area is of sufficient quality and disturbance by humans is minimal. Even in areas with a high density of humans, in the absence of directed persecution, species like Sarus Crane maintain territories as small as 5 ha when agricultural crops and landscape conditions are suitable. In contrast, red-crowned crane territories may require 500 hectares, and pairs may defend even larger territories than that, up to several thousand hectares. Territory defence is either acoustic with both birds performing the unison call, or more rarely, physical with attacks usually by the male. Because of this, females are much less likely to retain the territory than males in the event of the death of a partner. Rarely, breeding territorial crane pairs allow a third crane into the territory to form polygynous or polyandrous trios that improves the chances of survival of the pair's chicks. Trios of Sarus cranes were seen largely in marginal habitats and third birds were young suggesting that third cranes would benefit by gaining experience.
In mythology and symbolism
The cranes' beauty and spectacular mating dances have made them highly symbolic birds in many cultures with records dating back to ancient times. Crane mythology can be found in cultures around the world, from India to the Aegean, Arabia, China, Korea, Japan, Australia, and North America.
The Sanskrit epic poet Valmiki was inspired to write the first śloka couplet by the pathos of seeing a male sarus crane shot while dancing with its mate.
In Mecca, in pre-Islamic Arabia, Allāt, Uzza, and Manāt were believed to be the three chief goddesses of Mecca, they were called the "three exalted cranes" (, an obscure word on which 'crane' is the usual gloss). See The Satanic Verses for the best-known story regarding these three goddesses.
The Greek for crane is (), which gives us the cranesbill, or hardy geranium. The crane was a bird of omen. In the tale of Ibycus and the cranes, a thief attacked Ibycus (a poet of the sixth century BCE) and left him for dead. Ibycus called to a flock of passing cranes, which followed the attacker to a theater and hovered over him until, stricken with guilt, he confessed to the crime.
Pliny the Elder wrote that cranes would appoint one of their number to stand guard while they slept. The sentry would hold a stone in its claw, so that if it fell asleep, it would drop the stone and waken. A crane holding a stone in its claw is a well-known symbol in heraldry, and is known as a crane in its vigilance. Notably, however, the crest of Clan Cranstoun depicts a sleeping crane still in vigilance and holding the rock in its raised claw.
Aristotle describes the migration of cranes in the History of Animals, adding an account of their fights with Pygmies as they wintered near the source of the Nile. Battles between cranes and dwarf peoples, or geranomachy, is a widespread motif of antiquity and come from China and Arabia at least from the fifth century. Aristotle describes as untruthful an account that the crane carries a touchstone inside it that can be used to test for gold when vomited up.
Greek and Roman myths often portrayed the dance of cranes as a love of joy and a celebration of life, and the crane was often associated with both Apollo and Hephaestus.
In pre-modern Ottoman Empire, sultans would sometimes present a piece of crane feather (Turkish: ) to soldiers of any group in the army (janissaries, sipahis, etc.) who performed heroically during a battle. Soldiers would attach this feather to their caps or headgear which would give them some sort of a rank among their peers.
Throughout Asia, the crane is a symbol of happiness and eternal youth. In China, several styles of kung fu take inspiration from the movements of cranes in the wild, the most famous of these styles being Wing Chun, Hung Gar (tiger crane), and the Shaolin Five Animals style of fighting. Crane movements are well known for their fluidity and grace.
In Japan, the crane is one of the mystical or holy creatures (others include the dragon and the tortoise) and symbolizes good fortune and longevity because of its fabled life span of a thousand years. The crane is one of the subjects in the tradition of origami, or paper folding. An ancient Japanese legend promises that anyone who folds a thousand origami cranes will be granted a wish by a crane. In northern Hokkaidō, the women of the Ainu people performed a crane dance that was captured in 1908 in a photograph by Arnold Genthe. After World War II, the crane came to symbolize peace and the innocent victims of war through the story of schoolgirl Sadako Sasaki and her thousand origami cranes. Suffering from leukemia as a result of the atomic bombing of Hiroshima and knowing she was dying, she undertook to make a thousand origami cranes before her death at the age of 12. After her death, she became internationally recognised as a symbol of the innocent victims of war and remains a heroine to many Japanese girls.
| Biology and health sciences | Gruiformes | null |
198134 | https://en.wikipedia.org/wiki/Gruiformes | Gruiformes | The Gruiformes ( ) are an order containing a considerable number of living and extinct bird families, with a widespread geographical diversity. Gruiform means "crane-like".
Traditionally, a number of wading and terrestrial bird families that did not seem to belong to any other order were classified together as Gruiformes. These include 15 species of large cranes, about 145 species of smaller crakes and rails, as well as a variety of families comprising one to three species, such as the Heliornithidae, the limpkin, or the Psophiidae.
Other birds have been placed in this order more out of necessity to place them somewhere; this has caused the expanded Gruiformes to lack distinctive apomorphies. Recent studies indicate that these "odd Gruiformes" are if at all only loosely related to the cranes, rails, and relatives ("core Gruiformes").
Systematics
There are only two suprafamilial clades (natural groups) among the birds traditionally classified as Gruiformes. Rails (Rallidae), flufftails (Sarothruridae), finfoots and sungrebe (Heliornithidae), adzebills (Aptornithidae), trumpeters (Psophiidae), limpkin (Aramidae), and cranes (Gruidae) compose the suborder Grues and are termed "core-Gruiformes". These are the only true Gruiformes. The suborder Eurypygae includes the kagu (Rhynochetidae) and sunbittern (Eurypygidae). These are not even remotely related to Grues. The families of mesites or roatelos (Mesitornithidae), button-quails (Turnicidae), Australian plains-wanderer (Pedionomidae), seriemas (Cariamidae), and bustards (Otididae) each represent distinct and unrelated lineages. Many families known only from fossils have been assigned to the Gruiformes, e.g., Ergilornithidae, Phorusrhacidae, Messelornithidae, Eogruidae, Idiornithidae, Bathornithidae, to name just a few (see below). Though some of these are superficially 'crane-like' and the possibility exists that some may even be related to extant families traditionally included in the Gruiformes, there are no completely extinct families that can be confidently assigned to core-Gruiformes.
The traditional order Gruiformes was established by the influential German avian comparative anatomist Max Fürbringer (1888). Over the decades, many ornithologists suggested that members of the order were in fact more closely related to other groups (reviewed by Olson 1985, Sibley and Ahlquist 1990). For example, it was thought that sunbittern might be related to herons and that seriemas might be related to cuckoos. Olson and Steadman (1981) were first to correctly disband any of the traditional Gruiformes. They recognized that the Australian plains-wanderer (family Pedionomidae) was actually a member of the shorebirds (order Charadriiformes) based on skeletal characters. This was confirmed by Sibley and Ahlquist (1990) based on DNA–DNA hybridization and subsequently by Paton et al. (2003), Paton and Baker (2006) and Fain and Houde (2004, 2006). Sibley and Ahlquist furthermore removed button-quails (Turnicidae) from the Gruiformes based on large DNA–DNA hybridization distances to other supposed Gruiformes. However, it was not until the work of Paton et al. (2004) and Fain and Houde (2004, 2006) that the correct placement of buttonquails within the shorebirds (order Charadriiformes) was documented on the basis of phylogenetic analysis of multiple genetic loci. Using 12S ribosomal DNA sequences, Houde et al. (1997) were the first to present molecular genetic evidence of gruiform polyphyly, although apparently they were not convinced by it. However, on the basis of numerous additional sequence data, it has been shown decisively that the traditionally recognized Gruiformes consist of five to seven unrelated clades (Fain and Houde 2004, Ericson et al. 2006, Hackett et al. 2008).
Fain and Houde (2004) proposed that Neoaves are divisible into two clades, Metaves and Coronaves, although it has been suggested from the start that Metaves may be paraphyletic (Fain and Houde 2004, Ericson et al. 2006, Hackett et al. 2008). Sunbittern, kagu, and mesites all group within Metaves but all the other lineages of "Gruiformes" group either with a collection of waterbirds or landbirds within Coronaves. This division has been upheld by the combined analysis of as many as 30 independent loci (Ericson et al. 2006, Hackett et al. 2008), but is dependent on the inclusion of one or two specific loci in the analyses. One locus, i.e., mitochondrial DNA, contradicts the strict monophyly of Coronaves (Morgan-Richards et al. 2008), but phylogeny reconstruction based on mitochondrial DNA is complicated by the fact that few families have been studied, the sequences are heavily saturated (with back mutations) at deep levels of divergence, and they are plagued by strong base composition bias.
The kagu and sunbittern are one another's closest relatives. It had been proposed (Cracraft 2001) that they and the recently extinct adzebills (family Aptornithidae) from New Zealand constitute a distinct Gondwanan lineage. However, sunbittern and kagu are believed to have diverged from one another long after the break-up of Gondwanaland and the adzebills are in fact members of the Grues (Houde et al. 1997, Houde 2009). The seriemas and bustards represent distinct lineages within neoavian waterbirds.
Phylogeny
Gruiformes
Family †Songziidae Hou, 1990
Genus †Songzia Hou, 1990
Suborder Grui
Superfamily Gruoidea Vigors, 1825
Family †Geranoididae Wetmore, 1933
Family †Parvigruidae Mayr, 2005
Genus †Parvigrus Mayr, 2005
Genus †Rupelrallus Fischer, 1997
Family Aramidae Bonaparte, 1854 (limpkin)
Genus †Badistornis Wetmore, 1940
Genus Aramus Vieillot, 1816 [Courlili Buffon, 1781; Notherodius Wagler, 1827] (limpkins)
Family Psophiidae Bonaparte, 1831 (trumpeters)
Genus Psophia Linnaeus, 1758
Family †Eogruidae Wetmore, 1934
Genus †Sonogrus Kuročkin, 1981
Genus †Eogrus Wetmore, 1932 [Progrus Bendukidze, 1971]
Subfamily †Ergilornithinae
Genus †Ergilornis Kozlova, 1960
Genus †Amphipelargus Lydekker, 1891
Genus †Urmiornis Mecquenem, 1908
Family Gruidae (cranes)
Genus †Camusia Seguí, 2002
Subfamily Balearicinae Brasil, 1913
Genus †Aramornis Wetmore, 1926
Genus †Geranopsis Lydekker, 1871
Genus †Eobalearica Gureev, 1949
Genus Balearica Brisson, 1760 [Geranarchus Gloger, 1842] (crowned cranes)
Subfamily Gruinae Vigors, 1825
Genus †"Grus" conferta Miller & Sibley, 1942 [Olson & Rasmussen, 2001]
Genus †"Probalearica" mongolica Kurochkin, 1985
Genus †Palaeogrus Portis, 1885 [Palaeogrus Salvadori, 1884 nomen nudum]
Genus Antigone (Linnaeus, 1758)
Genus Leucogeranus (Pallas, 1773)
Genus Grus Brisson, 1760 non Moehring, 1758 [Anthropoides Vieillot, 1816; Bugeranus Gloger, 1841; Megalornis Gray, 1841; Leucogeranus Bonaparte, 1855; Mathewsena Iredale, 1914; Mathewsia Iredale, 1911; Limnogeranus Sharpe, 1893; Laomedontia Reichenbach, 1852; Philorchemon Gloger, 1842; Scops Gray, 1840 non Moehring, 1758 non Bruennich, 1772 npn Savigny, 1809] (cranes)
Suborder Ralli
Family †Aptornithidae (adzebills)
Genus †Aptornis
Family Sarothruridae (flufftails)
Genus Mentocrex Peters, 1933 (wood rails)
Genus Sarothrura Heine, 1890 non Hasselt, 1823 [Corethrura Reichenbach, 1849 non Hope, 1843 non Gray, 1846; Daseioura Penhallurick, 2003] (flufftails)
Family Heliornithidae Gray, 1841 (finfoots and sungrebe)
Genus Heliopais Sharpe, 1893 (Asian/masked finfoots)
Genus Podica Lesson, 1831 [Rhigelura Wagler, 1832; Podoa Bonaparte, 1857 non Illiger, 1811] (African finfoots)
Genus Heliornis Bonnaterre, 1791 [Podoa Illiger, 1811 non Bonaparte, 1857; Plotoides Brookes, 1830; Podia Swainson, 1837] (sungrebe, American finfoot)
Family Rallidae (crakes, moorhens, gallinules, and rails)
Genus †Aletornis Marsh, 1872 [Protogrus]
Genus †Australlus Worthy & Boles, 2011
Genus †Baselrallus De Pietri & Mayr, 2014
Genus †Belgirallus Mayr & Smith, 2001
Genus †Capellirallus Falla, 1954 (snipe-billed rail)
Genus †Creccoides Shufeldt, 1892
Genus †Eocrex Wetmore, 1931
Genus †Euryonotus Mercerat, 1897
Genus †Fulicaletornis Lambrecht, 1933
Genus †Hovacrex Brodkorb, 1965 (Hova gallinule)
Genus †Ibidopsis Lydekker, 1891
Genus †Latipons Harrison & Walker, 1979
Genus †Miofulica Lambrecht, 1933
Genus †Miorallus Lambrecht, 1933
Genus †Nesophalaris Brodkorb & Dawson, 1962
Genus †Palaeoaramides Lambrecht, 1933
Genus †Palaeorallus Wetmore, 1931
Genus †Paraortygometra Lambrecht, 1933
Genus †Parvirallus Harrison & Walker, 1979
Genus †Pastushkinia Zelenkov, 2013
Genus †Quercyrallus Lambrecht, 1933
Genus †Rallicrex Lambrecht, 1933
Genus †Rhenanorallus Mayr, 2010
Genus †Vitirallus Worthy, 2004 (Viti Levu rails)
Genus †Wanshuina Hou, 1994
Genus †Youngornis Yeh, 1981
Genus †Rallidae gen. et sp. indet. [Fulica podagrica (partim)] (Barbados rail)
Genus †Rallidae gen. et sp. indet. (Easter Island rail)
Genus †Rallidae gen. et sp. indet. (Fernando de Noronha rail)
Genus †Rallidae gen. et sp. indet. (Tahitian "goose")
Genus †Rallidae gen. et sp. indet. (Bokaak "bustard")
Genus †Rallidae gen. et sp. indet. ('Amsterdam Island' rail)
Genus Rougetius Bonaparte, 1856 (Rouget's Rails)
Subfamily Rallinae Rafinesque, 1815
Genus †Pleistorallus Worthy, 1997 (Fleming's rails)
Genus Anurolimnas Sharpe, 1893 (Chestnut-headed Crakes)
Genus Biensis (Madagascan Rails)
Genus Rallicula Schlegel, 1871 [Corethruropsis Salvadori, 1876] (forest-rails)
Genus Rallus Linnaeus, 1758 [†Epirallus Miller, 1942]
Genus †Aphanapteryx von Frauenfeld, 1868 [Pezocrex Hachisuka, 1953] (Mauritius/Red rails)
Genus †Erythromachus Milne-Edwards, 1873 (Rodriquez rails)
Genus Dryolimnas Sharpe, 1893
Genus Crex Bechstein, 1803 [Crecopsis Sharpe, 1893] (greater crakes)
Genus Lewinia Gray, 1855 [Aramidopsis Sharpe, 1893; Donacias Heine & Reichenow, 1890; Hyporallus Iredale & Mathews, 1926]
Genus Canirallus Bonaparte, 1856 (grey-throated rail)
Genus Gymnocrex Salvadori, 1875 (bare-faced rails)
Genus Gallirallus Lafresnaye, 1841 [Tricholimnas Sharpe, 1893; Nesoclopeus Peters, 1932; Cabalus Hutton, 1874; Habropteryx Stresemann, 1932; Eulabeornis Gould, 1844; †Diaphorapteryx Forbes, 1893; Hypotaenidia Reichenbach, 1853; Sylvestrornis Mathews, 1928]
Subfamily Gallinulinae Gray, 1840
Tribe Pardirallini Livezey, 1998 [Aramidinae] (Wood-rails & allies)
Genus Pardirallus Bonaparte, 1856 [Ortygonax Heine, 1890]
Genus Mustelirallus Bonaparte, 1858 [Neocrex Sclater & Salvin, 1869; Cyanolimnas Barbour & Peters, 1927]
Genus Amaurolimnas Sharpe 1893 (Rufous rails; Uniform crakes)
Genus Aramides Pucheran, 1845
Tribe Gallinulini Gray, 1840 [Fulicarinae (Nitzsch, 1820) sensu Livezey, 1998]
Genus Tribonyx Du Bus de Gisignies, 1840 [Brachyptrallus Lafresnaye, 1840; Microtribonyx Sharpe, 1893] (native-hens)
Genus Porzana Vieillot, 1816 [Limnobaenus Sundevall, 1872; Phalaridion Kaup, 1829; Porzanoidea Mathews, 1912; Porzanoides Condon, 1975; Rallites Pucheran, 1845; Schoenocrex Roberts, 1922; Porphyriops Pucheran, 1845]
Genus Paragallinula Sangster, García-R & Trewick, 2015 (Lesser Moorhen)
Genus Gallinula Brisson, 1760 [Hydrogallina Lacépède, 1799; Stagnicola Brehm, 1831; Porphyriornis Allen, 1892 Pareudiastes Hartlaub & Finsch, 1871 Edithornis]
Genus Fulica Linnaeus, 1758 [†Palaeolimnas Forbes, 1893]
Subfamily Porphyrioninae Reichenbach, 1849
Tribe Porphyrionini Reichenbach, 1849 (Purple gallinules & swamphens)
Genus †Aphanocrex Wetmore, 1963 (St. Helena swamphens)
Genus Porphyrio Brisson, 1760 [Notornis Owen, 1848]
Tribe Himantornithini Bonaparte, 1856 (Bush-hens & Waterhens)
Genus Himantornis Hartlaub, 1855 (Nkulenga rails)
Genus Megacrex D'Albertis & Salvadori, 1879 (New Guinea Flightless Rails)
Genus Aenigmatolimnas (Striped Crakes)
Genus Gallicrex Blyth, 1852 [Gallinulopha Bonaparte, 1854; Hypnodes Reichenbach, 1853] (Watercocks)
Genus Amaurornis Reichenbach, 1853 [Erythra Reichenbach, 1853; Pisynolimnas Heine & Reichenow, 1890; Poliolimnas Sharpe, 1893] (Bush-hen)
Tribe Zaporniini Des Murs, 1860 (Old world crakes)
Genus Rallina Gray, 1846 [Euryzona Gray, 1855; Tomirdus Mathews, 1912] (chestnut-rails)
Genus Zapornia Stephens, 1824 [Limnocorax Peters, 1854; Limnobaenus; Corethrura Grey, 1846]
Tribe Laterallini Tif, 2014 (New world crakes)
Genus Micropygia Bonaparte, 1856 (Ocellated Crakes)
Genus Rufirallus (russet-crowned crake)
Genus Laterallus Gray, 1855 (ruddy crakes)
Genus Coturnicops Gray, 1855 (barred-backed crakes)
Genus Hapalocrex (Yellow-breasted Crakes)
Genus Limnocrex
Genus Mundia Bourne, Ashmole & Simmons, 2003 (Ascension Island Crakes)
Genus Creciscus Cabanis, 1857 [Atlantisia Lowe, 1923] (blackish crakes)
Not placed in family
Genus †Nesotrochis Wetmore, 1918 (West Indian cave-rails)
When considered to be monophyletic, it was assumed that Gruiformes was among the more ancient of avian lineages. The divergence of "gruiforms" among "Metaves" and "Coronaves" is proposed to be the first divergence among Neoaves, far predating the Cretaceous–Paleogene extinction event c. 66 mya (Houde 2009). No unequivocal basal gruiforms are known from the fossil record. However, there are several genera that are not unequivocally assignable to the known families and that may occupy a more basal position:
Propelargus (Late Eocene/Early Oligocene of Quercy, France) – cariamid or idornithid
Rupelrallus (Early Oligocene of Germany) – rallid? parvigruid?
Badistornis (Brule Middle Oligocene of Shannon County, Missouri) – aramid?
Probalearica (Late Oligocene? – Middle Pliocene of Florida, France?, Moldavia and Mongolia) – gruid? A nomen dubium?
"Gruiformes" gen. et sp. indet. MNZ S42623 (Bathans Early/Middle Miocene of Otago, New Zealand) – Aptornithidae?
Aramornis (Sheep Creek Middle Miocene of Snake Creek Quarries, U.S.) – gruid? aramid?
Euryonotus (Pleistocene of Argentina) – rallid?
Other even more enigmatic fossil birds and five living families are occasionally suggested to belong into this order, such as the proposed Late Cretaceous family Laornithidae and the following taxa:
Family †Gastornithidae (diatrymas) (fossil)
Family †Messelornithidae (Messel-birds)
Family †Salmilidae (fossil) – distinct order (Cariamiformes)
Family †Geranoididae (fossil) – distinct order (Cariamiformes); however, Mayr (2016) argued they might be members of Gruiformes, specifically stem group representatives of the Gruoidea.
Family †Bathornithidae (fossil) – distinct order (Cariamiformes)
Family †Idiornithidae (fossil) – distinct order (Cariamiformes)
Family †Phorusrhacidae (terror birds) (fossil) – distinct order (Cariamiformes)
Family Cariamidae (seriemas) – Neoavian landbirds – distinct order (Cariamiformes)
Family Otididae (bustards) – Neoavian waterbirds – distinct order
Family Eurypygidae (sunbittern) – prospective "Metaves" – new order Eurypygiformes together with kagu
Family Rhynochetidae (kagu) – prospective "Metaves" – new order Eurypygiformes together with sunbittern
Family Mesitornithidae (mesites, roatelos, monias) prospective "Metaves" – distinct order
Family Turnicidae (buttonquails) moved to already existing order Charadriiformes together with plains wanderer
Family Pedionomidae (plains wanderer) moved to already existing order Charadriiformes together with buttonquails
Horezmavis (Bissekty Late Cretaceous of Kyzyl Kum, Uzbekistan)
Telmatornis (Navesink Late Cretaceous?)
Amitabha (Bridger middle Eocene of Forbidden City, Wyoming) – rallid?
Eobalearica (Ferghana Late? Eocene of Ferghana, Uzbekistan) – gruid?
"Phasianus" alfhildae (Washakie B Late Eocene of Haystack Butte, U.S.)
Talantatos (Late Eocene of Paris Bain, France)
Telecrex (Irdin Manha Late Eocene of Chimney Butte, China) – rallid?
Neornithes incerta sedis (Late Paleocene/Early Eocene of Ouled Abdoun Basin, Morocco)
Aminornis (Deseado Early Oligocene of Rio Deseado, Argentina) – aramid?
Loncornis (Deseado Early Oligocene of Rio Deseado, Argentina) – aramid?
Riacama (Deseado Early Oligocene of Argentina)
Smiliornis (Deseado Early Oligocene of Argentina)
Pseudolarus (Deseado Early Oligocene – Miocene of Argentina) – gruiform?
Gnotornis (Brule Late Oligocene of Shannon County, Missouri) – aramid?
Anisolornis (Santa Cruz Middle Miocene of Karaihen, Argentina) – aramid?
Occitaniavis – cariamid or idiornithid, includes Geranopsis elatus
| Biology and health sciences | Gruiformes | null |
198154 | https://en.wikipedia.org/wiki/Beak | Beak | The beak, bill, or rostrum is an external anatomical structure found mostly in birds, but also in turtles, non-avian dinosaurs and a few mammals. A beak is used for pecking, grasping, and holding (in probing for food, eating, manipulating and carrying objects, killing prey, or fighting), preening, courtship, and feeding young. The terms beak and rostrum are also used to refer to a similar mouth part in some ornithischians, pterosaurs, cetaceans, dicynodonts, rhynchosaurs, anuran tadpoles, monotremes (i.e. echidnas and platypuses, which have a bill-like structure), sirens, pufferfish, billfishes, and cephalopods.
Although beaks vary significantly in size, shape, color and texture, they share a similar underlying structure. Two bony projections–the upper and lower mandibles–are covered with a thin keratinized layer of epidermis known as the rhamphotheca. In most species, two holes called nares lead to the respiratory system.
Etymology
Although the word "beak" was, in the past, generally restricted to the sharpened bills of birds of prey, in modern ornithology, the terms beak and bill are generally considered to be synonymous. The word, which dates from the 13th century, comes from the Middle English bec (via Anglo French), which itself comes from the Latin beccus.
Anatomy
Although beaks vary significantly in size and shape from species to species, their underlying structures have a similar pattern. All beaks are composed of two jaws, generally known as the maxilla (upper) and mandible (lower). The upper, and in some cases the lower, mandibles are strengthened internally by a complex three-dimensional network of bony spicules (or trabeculae) seated in soft connective tissue and surrounded by the hard outer layers of the beak. The avian jaw apparatus is made up of two units: one four-bar linkage mechanism and one five-bar linkage mechanism.
Mandibles
The upper mandible is supported by a three-pronged bone called the intermaxillary. The upper prong of this bone is embedded into the forehead, while the two lower prongs attach to the sides of the skull. At the base of the upper mandible a thin sheet of nasal bones is attached to the skull at the nasofrontal hinge, which gives mobility to the upper mandible, allowing it to move upward and downward.
The base of the upper mandible, or the roof when seen from the mouth, is the palate; the palate's structure differs greatly in the ratites. Here, the vomer is large and connects with premaxillae and maxillopalatine bones in a condition termed as a "paleognathous palate". All other extant birds have a narrow forked vomer that does not connect with other bones and is then termed as neognathous. The shape of these bones varies across the bird families.
The lower mandible is supported by a bone known as the inferior maxillary bone—a compound bone composed of two distinct ossified pieces. These ossified plates (or rami), which can be U-shaped or V-shaped, join distally (the exact location of the joint depends on the species) but are separated proximally, attaching on either side of the head to the quadrate bone. The jaw muscles, which allow the bird to close its beak, attach to the proximal end of the lower mandible and to the bird's skull. The muscles that depress the lower mandible are usually weak, except in a few birds such as the starlings and the extinct huia, which have well-developed digastric muscles that aid in foraging by prying or gaping actions. In most birds, these muscles are relatively small as compared to the jaw muscles of similarly sized mammals.
Rhamphotheca
The outer surface of the beak consists of a thin sheath of keratin called the rhamphotheca, which can be subdivided into the rhinotheca of the upper mandible and the gnathotheca of the lower mandible. The covering arises from the Malpighian layer of the bird's epidermis,
growing from plates at the base of each mandible. There is a vascular layer between the rhamphotheca and the deeper layers of the dermis, which is attached directly to the periosteum of the bones of the beak. The rhamphotheca grows continuously in most birds, and in some species, the color varies seasonally.
In some alcids, such as the puffins, parts of the rhamphotheca are shed each year after the breeding season, while some pelicans shed a part of the bill called a "bill horn" that develops in the breeding season.
While most extant birds have a single seamless rhamphotheca, species in a few families, including the albatrosses and the emu, have compound rhamphothecae that consist of several pieces separated and defined by softer keratinous grooves.
Studies have shown that this was the primitive ancestral state of the rhamphotheca, and that the modern simple rhamphotheca resulted from the gradual loss of the defining grooves through evolution.
Tomia
The tomia (singular tomium) are the cutting edges of the two mandibles. In most birds, they range from being rounded to slightly sharp, but some species have evolved structural modifications that allow them to handle their typical food sources better.
Granivorous (seed-eating) birds, for example, have ridges in their tomia, which help the bird to slice through a seed's outer hull.
Most falcons have a sharp projection along the upper mandible, with a corresponding notch on the lower mandible. They use this "tooth" to sever their prey's vertebrae fatally or to rip insects apart. Some kites, principally those that prey on insects or lizards, also have one or more of these sharp projections,
as do the shrikes.
The tomial teeth of falcons are underlain by bone, while the shrike tomial teeth are entirely keratinous. Some fish-eating species, e.g., the mergansers, have sawtooth serrations along their tomia, which help them to keep hold of their slippery, wriggling prey.
Birds in roughly 30 families have tomia lined with tight bunches of very short bristles along their entire length. Most of these species are either insectivores (preferring hard-shelled prey) or snail eaters, and the brush-like projections may help to increase the coefficient of friction between the mandibles, thereby improving the bird's ability to hold hard prey items.
Serrations on hummingbird bills, found in 23% of all hummingbird genera, may perform a similar function, allowing the birds to effectively hold insect prey. They may also allow shorter-billed hummingbirds to function as nectar thieves, as they can more effectively hold and cut through long or waxy flower corollas.
In some cases, the color of a bird's tomia can help to distinguish between similar species. The snow goose, for example, has a reddish-pink bill with black tomia, while the whole beak of the similar Ross's goose is pinkish-red, without darker tomia.
Culmen
The culmen is the dorsal ridge of the upper mandible. Likened by ornithologist E. Coues to the ridge line of a roof, it is the "highest middle lengthwise line of the bill" and runs from the point where the upper mandible emerges from the forehead's feathers to its tip. The bill's length along the culmen is one of the regular measurements made during bird banding (ringing)
and is particularly useful in feeding studies. There are several standard measurements which can be made—from the beak's tip to the point where feathering starts on the forehead, from the tip to the anterior edge of the nostrils, from the tip to the base of the skull, or from the tip to the cere (for raptors and owls)—and scientists from various parts of the world generally favor one method over another. In all cases, these are chord measurements (measured in a straight line from point to point, ignoring any curve in the culmen) taken with calipers.
The shape or color of the culmen can also help with the identification of birds in the field. For example, the culmen of the parrot crossbill is strongly decurved, while that of the very similar red crossbill is more moderately curved.
The culmen of a juvenile common loon is all dark, while that of the very similarly plumaged juvenile yellow-billed loon is pale towards the tip.
Gonys
The gonys is the ventral ridge of the lower mandible, created by the junction of the bone's two rami, or lateral plates. The proximal end of that junction—where the two plates separate—is known as the gonydeal angle or gonydeal expansion. In some gull species, the plates expand slightly at that point, creating a noticeable bulge; the size and shape of the gonydeal angle can be useful in identifying between otherwise similar species. Adults of many species of large gulls have a reddish or orangish gonydeal spot near the gonydeal expansion. This spot triggers begging behavior in gull chicks. The chick pecks at the spot on its parent's bill, which in turn stimulates the parent to regurgitate food.
Commissure
Depending on its use, commissure may refer to the junction of the upper and lower mandibles, or alternately, to the full-length apposition of the closed mandibles, from the corners of the mouth to the tip of the beak.
Gape
In bird anatomy, the gape is the interior of the open mouth of a bird, and the gape flange is the region where the two mandibles join together at the base of the beak.
The width of the gape can be a factor in the choice of food.
Gapes of juvenile altricial birds are often brightly coloured, sometimes with contrasting spots or other patterns, and these are believed to be an indication of their health, fitness and competitive ability. Based on that, the parents decide how to distribute food among the chicks in the nest.
Some species, especially in the families Viduidae and Estrildidae, have bright spots on the gape known as gape tubercles or gape papillae. These nodular spots are conspicuous even in low light.
A study examining the nestling gapes of eight passerine species found that the gapes were conspicuous in the ultraviolet spectrum (visible to birds but not to humans).
Parents may, however, not rely solely on the gape coloration, and other factors influencing their decision remain unknown.
Red gape color has been shown in several experiments to induce feeding. An experiment in manipulating brood size and immune system with barn swallow nestlings showed the vividness of the gape was positively correlated with T-cell–mediated immunocompetence, and that larger brood size and injection with an antigen led to a less vivid gape.
Conversely, the red gape of the common cuckoo (Cuculus canorus) did not induce extra feeding in host parents.
Some brood parasites, such as the Hodgson's hawk-cuckoo (C. fugax), have colored patches on the wing that mimic the gape color of the parasitized species.
When born, the chick's gape flanges are fleshy. As it grows into a fledgling, the gape flanges remain somewhat swollen and can thus be used to recognize that a particular bird is young.
By the time it reaches adulthood, the gape flanges will no longer be visible.
Nares
Most species of birds have external nares (nostrils) located somewhere on their beak. The nares are two holes—circular, oval or slit-like in shape—which lead to the nasal cavities within the bird's skull, and thus to the rest of the respiratory system. In most bird species, the nares are located in the basal third of the upper mandible. Kiwis are a notable exception; their nares are located at the tip of their bills. A handful of species have no external nares. Cormorants and darters have primitive external nares as nestlings, but these close soon after the birds fledge; adults of these species (and gannets and boobies of all ages, which also lack external nostrils) breathe through their mouths. There is typically a septum made of bone or cartilage that separates the two nares, but in some families (including gulls, cranes, and New World vultures), the septum is missing. While the nares are uncovered in most species, they are covered with feathers in a few groups of birds, including grouse and ptarmigans, crows, and some woodpeckers. The feathers over a ptarmigan's nostrils help to warm the air it inhales,
while those over a woodpecker's nares help to keep wood particles from clogging its nasal passages.
Species in the bird order Procellariformes have nostrils enclosed in double tubes which sit atop or along the sides of the upper mandible. These species, which include the albatrosses, petrels, diving petrels, storm petrels, fulmars and shearwaters, are widely known as "tubenoses".
A number of species, including the falcons, have a small bony tubercule which projects from their nares. The function of the tubercule is unknown. Some scientists suggest it acts as a baffle, slowing down or diffusing airflow into the nares (thus allowing the bird to continue breathing without damaging its respiratory system) during high-speed dives, but this theory has not been proved experimentally. Not all species which fly at high speeds have these tubercules, while some species which fly at low speeds have them.
Operculum
The nares of some birds are covered by an operculum (plural opercula), a membraneous, horny or cartilaginous flap.
In diving birds, the operculum keeps water out of the nasal cavity; when the birds dive, the impact force of the water closes the operculum.
Some species which feed on flowers have opercula to help to keep pollen from clogging their nasal passages, while the opercula of the two species of Attagis seedsnipe help to keep dust out.
The nares of nestling tawny frogmouths are covered with large dome-shaped opercula, which help to reduce the rapid evaporation of water vapor, and may also help to increase condensation within the nostrils themselves—both critical functions, since the nestlings get fluids only from the food their parents bring them. The opercula shrink as the birds age, disappearing completely by the time they reach adulthood.
In pigeons, the operculum has evolved into a soft swollen mass that sits at the base of the bill, above the nares; though it is sometimes referred to as the cere, this is a different structure. Tapaculos are the only birds known to have the ability to move their opercula.
Rosette
Some species like the puffin, have a fleshy rosette, sometimes called a "gape rosette"
at the corners of the beak. In the puffin, it is grown as part of its display plumage.
Cere
Birds from a handful of families—including raptors, owls, skuas, parrots, pigeons, turkeys and curassows—have a waxy structure called a cere (from the Latin cera, which means "wax") or ceroma
which covers the base of their bill. This structure typically contains the nares, except in the owls, where the nares are distal to the cere. Although it is sometimes feathered in parrots,
the cere is typically bare and often brightly colored. In raptors, the cere is a sexual signal which indicates the "quality" of a bird; the orangeness of a Montagu's harrier's cere, for example, correlates to its body mass and physical condition.
The cere color of young Eurasian scops-owls has an ultraviolet (UV) component, with a UV peak that correlates to the bird's mass. A chick with a lower body mass has a UV peak at a higher wavelength than a chick with a higher body mass does. Studies have shown that parent owls preferentially feed chicks with ceres that show higher wavelength UV peaks, that is, lighter-weight chicks.
The color or appearance of the cere can be used to distinguish between males and females in some species. For example, the male great curassow has a yellow cere, which the female (and young males) lack.
The male budgerigar's cere is royal blue, while the female's is a very pale blue, white, or brown.
Nail
All birds of the family Anatidae (ducks, geese, and swans) have a nail, a plate of hard horny tissue at the tip of the beak. The shield-shaped structure, which sometimes spans the entire width of the beak, is often bent at the tip to form a hook.
It serves different purposes depending on the bird's primary food source. Most species use their nails to dig seeds out of mud or vegetation,
while diving ducks use theirs to pry molluscs from rocks.
There is evidence that the nail may help a bird to grasp objects. Species which use strong grasping motions to secure their food (such as when catching and holding onto a large squirming frog) have very wide nails.
Certain types of mechanoreceptors, nerve cells that are sensitive to pressure, vibration, or touch, are located under the nail.
The shape or color of the nail can sometimes be used to help distinguish between similar-looking species or between various ages of waterfowl. For example, the greater scaup has a wider black nail than does the very similar lesser scaup.
Juvenile "grey geese" have dark nails, while most adults have pale nails.
The nail gave the wildfowl family one of its former names: "Unguirostres" comes from the Latin ungus, meaning "nail" and rostrum, meaning "beak".
Rictal bristles
Rictal bristles are stiff hair-like feathers that arise around the base of the beak.
They are common among insectivorous birds, but are also found in some non-insectivorous species.
Their function is uncertain, although several possibilities have been proposed. They may function as a "net", helping in the capture of flying prey, although to date, there has been no empirical evidence to support this idea.
There is some experimental evidence to suggest that they may prevent particles from striking the eyes if, for example, a prey item is missed or broken apart on contact. They may also help to protect the eyes from particles encountered in flight, or from casual contact from vegetation. There is also evidence that the rictal bristles of some species may function tactilely, in a manner similar to that of mammalian whiskers (vibrissae). Studies have shown that Herbst corpuscles, mechanoreceptors sensitive to pressure and vibration, are found in association with rictal bristles. They may help with prey detection, with navigation in darkened nest cavities, with the gathering of information during flight or with prey handling.
Egg tooth
Full-term chicks of most bird species have a small sharp, calcified projection on their beak, which they use to chip their way out of their egg.
Commonly known as an egg tooth, this white spike is generally near the tip of the upper mandible, though some species have one near the tip of their lower mandible instead, and a few species have one on each mandible.
Despite its name, the projection is not an actual tooth, as the similarly-named projections of some reptiles are; instead, it is part of the integumentary system, as are claws and scales.
The hatching chick first uses its egg tooth to break the membrane around an air chamber at the wide end of the egg. Then it pecks at the eggshell while turning slowly within the egg, eventually (over a period of hours or days) creating a series of small circular fractures in the shell.
Once it has breached the egg's surface, the chick continues to chip at it until it has made a large hole. The weakened egg eventually shatters under the pressure of the bird's movements.
The egg tooth is so critical to a successful escape from the egg that chicks of most species will perish unhatched if they fail to develop one. However, there are a few species which do not have egg teeth. Megapode chicks have an egg tooth while still in the egg but lose it before hatching, while kiwi chicks never develop one; chicks of both families escape their eggs by kicking their way out.
Most chicks lose their egg teeth within a few days of hatching, although petrels keep theirs for nearly three weeks and marbled murrelets have theirs for up to a month.
Generally, the egg tooth drops off, though in songbirds it is resorbed.
Color
The color of a bird's beak results from concentrations of pigments—primarily melanins and carotenoids—in the epidermal layers, including the rhamphotheca.
Eumelanin, which is found in the bare parts of many bird species, is responsible for all shades of gray and black; the denser the deposits of pigment found in the epidermis, the darker the resulting color. Phaeomelanin produces "earth tones" ranging from gold and rufous to various shades of brown. Although it is thought to occur in combination with eumelanin in beaks which are buff, tan, or horn-colored, researchers have yet to isolate phaeomelanin from any beak structure. More than a dozen types of carotenoids are responsible for the coloration of most red, orange, and yellow beaks.
The hue of the color is determined by the precise mix of red and yellow pigments, while the saturation is determined by the density of the deposited pigments. For example, bright red is created by dense deposits of mostly red pigments, while dull yellow is created by diffuse deposits of mostly yellow pigments. Bright orange is created by dense deposits of both red and yellow pigments, in roughly equal concentrations. Beak coloration helps to make displays using those beaks more obvious. In general, beak color depends on a combination of the bird's hormonal state and diet. Colors are typically brightest as the breeding season approaches, and palest after breeding.
Birds are capable of seeing colors in the ultraviolet range, and some species are known to have ultraviolet peaks of reflectance (indicating the presence of ultraviolet color) on their beaks. The presence and intensity of these peaks may indicate a bird's fitness, sexual maturity or pair bond status.
King and emperor penguins, for example, show spots of ultraviolet reflectance only as adults. These spots are brighter on paired birds than on courting birds. The position of such spots on the beak may be important in allowing birds to identify conspecifics. For instance, the very similarly-plumaged king and emperor penguins have UV-reflective spots in different positions on their beaks.
Dimorphism
The size and shape of the beak can vary across species as well as between them; in some species, the size and proportions of the beak vary between males and females. That allows the sexes to utilize different ecological niches, thereby reducing intraspecific competition.
For example, females of nearly all shorebirds have longer bills than males of the same species,
and female American avocets have beaks which are slightly more upturned than those of males.
Males of the larger gull species have bigger, stouter beaks than those of females of the same species, and immatures can have smaller, more slender beaks than those of adults.
Many hornbills show sexual dimorphism in the size and shape of both beaks and casques, and the female huia's slim, decurved bill was nearly twice as long as the male's straight, thicker one.
Color can also differ between sexes or ages within a species. Typically, such a color difference is due to the presence of androgens. For example, in house sparrows, melanins are produced only in the presence of testosterone; castrated male house sparrows—like female house sparrows—have brown beaks. Castration also prevents the normal seasonal color change in the beaks of male black-headed gulls and indigo buntings.
Development
The beak of modern birds has a fused premaxillary bone, which is modulated by the expression of Fgf8 gene in the frontonasal ectodermal zone during embryonic development.
The shape of the beak is determined by two modules: the prenasal cartilage during early embryonic stage and the premaxillary bone during later stages. Development of the prenasal cartilage is regulated by genes Bmp4 and CaM, while that of the premaxillary bone is controlled by TGFβllr, β-catenin, and Dickkopf-3. TGFβllr codes for a serine/threonine protein kinase that regulates gene transcription upon ligand binding; previous work has highlighted its role in mammalian craniofacial skeletal development. β-catenin is involved in the differentiation of terminal bone cells. Dickkopf-3 codes for a secreted protein also known to be expressed in mammalian craniofacial development. The combination of these signals determines beak growth along the length, depth, and width axes. Reduced expression of TGFβllr significantly decreased the depth and length of chicken embryonic beak due to the underdevelopment of the premaxillary bone. Contrarily, an increase in Bmp4 signaling would result in a reduced premaxillary bone due to the overdevelopment of the prenasal cartilage, which takes up more mesenchymal cells for cartilage, instead of bone, formation.
Functions
Eating
Different species' beaks have evolved according to their diet; for example, raptors have sharp-pointed beaks that facilitate dissection and biting off of prey animals' tissue, whereas passerine birds that specialize in eating seeds with especially tough shells (such as grosbeaks and cardinals) have large, stout beaks with high compressive power (on the same principle as a human-devised nutcracker). Birds which fish for a living have beaks adapted for that pursuit; for example, pelicans' beaks are well adapted for scooping up and swallowing fish whole. Woodpeckers have beaks well adapted for pecking apart wood while hunting for their meals of arthropods.
Self-defensive pecking
Birds may bite or stab with their beaks to defend themselves.
Displays (for courtship, territoriality, or deterrence)
Some species use their beaks in displays of various sorts. As part of his courtship, for example, the male garganey touches his beak to the blue speculum feathers on his wings in a fake preening display, and the male mandarin duck does the same with his orange sail feathers. A number of species use a gaping, open beak in their fear and/or threat displays. Some augment the display by hissing or breathing heavily, while others clap their beak. Red-bellied woodpeckers at bird feeders often wave their formidable beaks at competing birds who get too close, clearly signaling "this seed's mine, you can't have it."
Sensory detection
The platypus uses its bill to navigate underwater, detect food, and dig. The bill contains electroreceptors and mechanoreceptors, causing muscular contractions to help detect prey. It is one of the few species of mammals to use electroreception.
The beaks of Kiwis, Ibises, and sandpipers have sensory pits in their beaks that allow them to sense vibrations.
The beaks of aquatic birds contain Grandry corpuscles, which assist in velocity detection while filter feeding.
Preening
The beak of birds plays a role in removing skin parasites (ectoparasites) such as lice. It is mainly the tip of the beak that does this. Studies have shown that inserting a bit to stop birds from using the tip results in increased parasite loads in pigeons.
Birds that have naturally deformed beaks have also been noted to have higher levels of parasites.
It is thought that the overhang at the end of the top portion of the beak (that is the portion that begins to curve downwards) slides against the lower beak to crush parasites.
This overhang of the beak is thought to be under stabilising natural selection. Very long beaks are thought to be selected against because they are prone to a higher number of breaks, as has been demonstrated in rock pigeons. Beaks with no overhang would be unable to effectively remove and kill ectoparasites as mentioned above. Studies have supported there is a selection pressure for an intermediate amount of overhang. Western scrub jays who had more symmetrical bills (i.e. those with less of an overhang), were found to have higher amounts of lice when tested. The same pattern has been found in surveys of Peruvian birds.
Additionally because of the role beaks play in preening, this is evidence for coevolution of the beak overhang morphology and body morphology of parasites. Artificially removing the ability to preen in birds, followed by readdition of preening ability was shown to result in changes in body size in lice. Once the ability of the birds to preen was reintroduced, the lice were found to show declines in body size suggesting they may evolve in response to preening pressures from birds who could respond in turn with changes in beak morphology.
Communicative percussion
A number of species, including storks, some owls, frogmouths and the noisy miner, use bill clapping as a form of communication. Some woodpecker species are known to use percussion as a courtship activity, whereas males get the (aural) attention of females from a distance and then impress them with the sound volume and pattern. This explains why humans are sometimes inconvenienced by pecking that clearly has no feeding purpose (such as when the bird pecks on sheet metal repeatedly).
Heat exchange
Studies have shown that some birds use their beaks to rid themselves of excess heat. The toco toucan, which has the largest beak relative to the size of its body of any bird species, is capable of modifying the blood flow to its beak. This process allows the beak to work as a "transient thermal radiator", reportedly rivaling an elephant's ears in its ability to radiate body heat.
Measurements of the bill sizes of several species of American sparrows found in salt marshes along the North American coastlines show a strong correlation with summer temperatures recorded in the locations where the sparrows breed; latitude alone showed a much weaker correlation. By dumping excess heat through their bills, the sparrows are able to avoid the water loss which would be required by evaporative cooling—an important benefit in a windy habitat where freshwater is scarce.
Several ratites, including the common ostrich, the emu and the southern cassowary, use various bare parts of their bodies (including their beaks) to dissipate as much as 40% of their metabolic heat production.
Alternately, studies have shown that birds from colder climates (higher altitudes or latitudes and lower environmental temperatures) have smaller beaks, lessening heat loss from that structure.
Billing
During courtship, mated pairs of many bird species touch or clasp each other's bills. Termed billing (also nebbing in British English),
this behavior appears to strengthen pair bonding. The amount of contact involved varies among species. Some gently touch only a part of their partner's beak while others clash their beaks vigorously together.
Gannets raise their bills high and repeatedly clatter them, the male puffin nibbles at the female's beak, the male waxwing puts his bill in the female's mouth and ravens hold each other's beaks in a prolonged "kiss".
Billing can also be used as a gesture of appeasement or subordination. Subordinate Canada jay routinely bill more dominant birds, lowering their body and quivering their wings in the manner of a young bird food begging as they do so.
A number of parasites, including rhinonyssids and Trichomonas gallinae are known to be transferred between birds during episodes of billing.
Use of the term extends beyond avian behavior; "billing and cooing" in reference to human courtship (particularly kissing) has been in use since Shakespeare's time, and derives from the courtship of doves.
Beak trimming
Because the beak is a sensitive organ with many sensory receptors, beak trimming (sometimes referred to as 'debeaking') is "acutely painful" to the birds it is performed on. It is nonetheless routinely done to intensively farmed poultry flocks, particularly laying and broiler breeder flocks, because it helps reduce the damage the flocks inflict on themselves due to a number of stress-induced behaviors, including cannibalism, vent pecking and feather pecking. A cauterizing blade or infrared beam is used to cut off about half of the upper beak and about a third of the lower beak. Pain and sensitivity can persist for weeks or months after the procedure, and neuromas can form along the cut edges. Food intake typically decreases for some period after the beak is trimmed. However, studies show that trimmed poultry's adrenal glands weigh less, and their plasma corticosterone levels are lower than those found in untrimmed poultry, indicating that they are less stressed overall.
A similar but separate practice, usually performed by an avian veterinarian or an experienced birdkeeper, involves clipping, filing or sanding the beaks of captive birds for health purposes–in order to correct or temporarily alleviate overgrowths or deformities and better allow the bird to go about its normal feeding and preening activities.
Among raptor keepers, the practice is commonly known as "coping".
Bill tip organ
The bill tip organ is a region found near the tip of the bill in several types of birds that forage particularly by probing. The region has a high density of nerve endings known as the corpuscles of Herbst. This consists of pits in the bill surface which in the living bird is occupied by cells that sense pressure changes. The assumption is that this allows the bird to perform 'remote touch', which means that it can detect movements of animals which the bird does not directly touch. Bird species known to have a 'bill-tip organ' include ibises, shorebirds of the family Scolopacidae, and kiwis.
There is a suggestion that across these species, the bill tip organ is better-developed among species foraging in wet habitats (water column, or soft mud) than in species using a more terrestrial foraging. However, it has been described in terrestrial birds too, including parrots, who are known for their dextrous extractive foraging techniques. Unlike probing foragers, the tactile pits in parrots are embedded in the hard keratin (or rhamphotheca) of the bill, rather than the bone, and along the inner edges of the curved bill, rather than being on the outside of the bill.
| Biology and health sciences | External anatomy and regions of the body | Biology |
198204 | https://en.wikipedia.org/wiki/Tetragonal%20crystal%20system | Tetragonal crystal system | In crystallography, the tetragonal crystal system is one of the 7 crystal systems. Tetragonal crystal lattices result from stretching a cubic lattice along one of its lattice vectors, so that the cube becomes a rectangular prism with a square base (a by a) and height (c, which is different from a).
Bravais lattices
There are two tetragonal Bravais lattices: the primitive tetragonal and the body-centered tetragonal.
The body-centered tetragonal lattice is equivalent to the primitive tetragonal lattice with a smaller unit cell, while the face-centered tetragonal lattice is equivalent to the body-centered tetragonal lattice with a smaller unit cell.
Crystal classes
The point groups that fall under this crystal system are listed below, followed by their representations in international notation, Schoenflies notation, orbifold notation, Coxeter notation and mineral examples.
In two dimensions
There is only one tetragonal Bravais lattice in two dimensions: the square lattice.
| Physical sciences | Crystallography | Physics |
198207 | https://en.wikipedia.org/wiki/Rec.%20601 | Rec. 601 | ITU-R Recommendation BT.601, more commonly known by the abbreviations Rec. 601 or BT.601 (or its former name CCIR 601), is a standard originally issued in 1982 by the CCIR (an organization, which has since been renamed as the International Telecommunication Union Radiocommunication sector) for encoding interlaced analog video signals in digital video form. It includes methods of encoding 525-line 60 Hz and 625-line 50 Hz signals, both with an active region covering 720 luminance samples and 360 chrominance samples per line. The color encoding system is known as YCbCr 4:2:2.
The Rec. 601 video raster format has been re-used in a number of later standards, including the ISO/IEC MPEG and ITU-T H.26x compressed formats, although compressed formats for consumer applications usually use chroma subsampling reduced from the 4:2:2 sampling specified in Rec. 601 to 4:2:0.
The standard has been revised several times in its history. Its seventh edition, referred to as BT.601-7, was approved in March 2011 and was formally published in October 2011.
Background and history
In the early 1980s, digital television equipment was beginning to emerge, but each manufacturer was developing their own proprietary digital versions of existing analog standards like PAL, SECAM, and NTSC.
At an ITU meeting in autumn 1981, CCIR Study Group 11 approved document 11/1027 describing the parameter values for a unified digital video format. This was adopted as Draft Rec. AA/11 "Encoding Parameters for Digital Television for Studios" by the CCIR Plenary Assembly in February 1982, later becoming ITU-R Rec. 601.
The key feature allowing a globally accepted digital standard were the use of component coding and choosing a luminance sampling frequency that was a common multiple of the line frequencies used in analog standards. This "orthogonal" sampling approach originated from Stanley Baron of NBC.
Preparation preceded the CCIR approval, including laboratory testing around the world to validate the proposed parameter values. International negotiations and efforts to build consensus were led by figures like Mark Krivosheev, Richard Green, and representatives from Japan and Europe.
Signal format
The Rec. 601 signal can be regarded as if it is a digitally encoded analog component video signal, and thus the sampling includes data for the horizontal and vertical sync and blanking intervals. Regardless of the frame rate, the luminance sampling frequency is 13.5 MHz. The samples are uniformly quantized using 8- or 10-bit PCM codes in the YCbCr domain.
For each 8-bit luminance sample, the nominal value to represent black is 16, and the value for white is 235. Eight-bit code values from 1 through 15 provide footroom and can be used to accommodate transient signal content such as filter undershoots. Similarly, code values 236 through 254 provide headroom and can be used to accommodate transient signal content such as filter overshoots. The values 0 and 255 are used to encode the sync pulses and are forbidden within the visible picture area. The Cb and Cr samples are unsigned and use the value 128 to encode the neutral color difference value, as used when encoding a white, grey or black area.
Primary chromaticities
Slightly different primaries are specified for the 625-line (PAL and SECAM) and 525-line (NTSC SMPTE C primaries) systems. Earlier versions of the standard (prior to BT.601-6, approved in January 2007) did not contain an explicit definition of the color primaries.
Transfer characteristics
Rec. 601 defines a nonlinear transfer function which is linear near 0 and then transfers to a gamma curve for the rest of the luminance range:
Awards
The CCIR received a 1982–83 Technology and Engineering Emmy Award for its development of the Rec. 601 standard.
| Physical sciences | Basics | Physics |
198290 | https://en.wikipedia.org/wiki/Chinchilla | Chinchilla | Chinchilla refers to either of two species (Chinchilla chinchilla and Chinchilla lanigera) of crepuscular rodents of the parvorder Caviomorpha, and are native to the Andes mountains in South America. They live in colonies called "herds" at high elevations up to . Historically, chinchillas lived in an area that included parts of Bolivia, Peru and Chile, but today, colonies in the wild are known only in Chile. Along with their relatives, viscachas, they make up the family Chinchillidae. They are also related to the chinchilla rat.
The chinchilla has the densest fur of all extant terrestrial mammals, with around 20,000 hairs per square centimeter and 50 hairs growing from each follicle. The chinchilla is named after the Chincha people of the Andes, who once wore its dense, velvet-like fur and ate their meat. By the end of the 19th century, chinchillas had become quite rare after being hunted for their notably soft fur. Most chinchillas currently used by the fur industry for clothing and other accessories are farm-raised. Domestic chinchillas descended from C. lanigera are sometimes kept as pets, and may be considered a type of pocket pet.
Species
The two living species of chinchilla are Chinchilla chinchilla (formerly known as Chinchilla brevicaudata) and Chinchilla lanigera. C. chinchilla has a shorter tail, a thicker neck and shoulders, and shorter ears than C. lanigera. The former species is currently facing extinction; the latter, though rare, can be found in the wild. Domesticated chinchillas are thought to be of the C. lanigera species.
Distribution and habitat
Chinchillas formerly occupied the coastal regions, hills, and mountains of Chile, Peru, Argentina, and Bolivia. Overexploitation caused the downturn of these populations and, as early as 1914, one scientist claimed that the species was headed for extinction. Five years of fieldwork (published in 2007) in Jujuy Province, Argentina, failed to find a single specimen. Populations in Chile were thought extinct by 1953, but the animal was found to inhabit an area in the Antofagasta Region in the late 1900s and early 2000s. The animal may be extinct in Bolivia and Peru, though one specimen found (in a restaurant in Cerro de Pasco) may hail from a native population.
In their native habitats, chinchillas live in burrows or crevices in rocks. They are agile jumpers and can jump up to . Predators in the wild include birds of prey, skunks, felines, snakes and canines. Chinchillas have a variety of defensive tactics, including spraying urine and releasing fur if bitten. In the wild, chinchillas have been observed eating plant leaves, fruits, seeds, and small insects.
In nature, chinchillas live in social groups that resemble colonies, but are properly called herds. Herd sizes can range from 14 members up to 100, and herding behavior is thought to promote both social interaction and protection from predators. They can breed any time of the year, though breeding season typically falls between May and November. They are typically monogamous. Their gestation period is 111 days, longer than most rodents. Due to this long pregnancy, chinchillas are born fully furred and with eyes open. Litters are usually small in number, predominantly two.
Conservation
Both species of chinchilla are currently listed as Endangered by the IUCN Red List of Threatened Species due to a severe population loss approximated at a 90% global population loss since 2001. The severe population decline has been caused by chinchilla hunting by humans. The long tailed-species was listed on the IUCN Red List as "Very rare and believed to be decreasing in numbers" in 1965. From 1982 to 1996, both species were listed as Indeterminate. In 2006, the long-tailed species was listed as "Vulnerable" while the short-tailed species was listed as "Critically Endangered". By 2008, both were listed as "Critically Endangered", and in 2016 they were reclassified as "Endangered" due to limited recovery in some areas.
Relationship with humans
Fur industry
Chinchilla fur trade on an international level goes back to the 16th century. Their fur is popular due to its extremely soft feel, which is caused by the sprouting of 25 hairs (on average) from each hair follicle. The color is usually very even, which makes it ideal for small garments or the lining of larger ones, though some large pieces can be made entirely from the fur. A single, full-length coat made from chinchilla fur may require as many as 150 pelts, as chinchillas are relatively small. Their use for fur led to the near extinction of one species(C.chinchilla), and put serious pressure on the other(C. lanigera). Though it is illegal to hunt wild chinchillas, they are now on the verge of becoming extinct because of continued poaching. Domesticated chinchillas are still bred for fur.
As pets
The domestic chinchilla is descended from Chinchilla lanigera, the long-tailed Chinchilla. They are the more common one in the wild, as the other species, Chinchilla chinchilla, or short-tailed Chinchilla, has been hunted nearly to extinction. Therefore, domestic chinchillas have thinner bodies, longer tails and larger ears. In the wild, the average life-span of a chinchilla is ten years; however, they could live up to 20 years in human care.
Chinchillas are popular pets, though they require extensive exercise and dental care, due to their teeth continually growing throughout their life span, and since they lack the ability to sweat, they require a temperature-controlled environment.
The animals instinctively clean their fur by taking dust baths, in which they roll around in special dust made of fine pumice, a few times a week; they do not bathe in water. Their thick fur resists parasites, such as fleas, and reduces loose dander.
Pet chinchillas require easy access to food, water, and hiding places, where they can sleep undisturbed for extended periods of time. Chinchillas are typically highly social creatures, so owners should interact often with their pets. They also have sensitive hearing and are easily startled by loud, unexpected noises.
In scientific research
Chinchillas have been used in research since the 1950s. Since the 1970s, the prime interest in chinchillas by researchers is their auditory system. Other research fields in which chinchillas are used as an animal model include the study of Chagas disease, gastrointestinal diseases, pneumonia, and listeriosis, as well as of Yersinia and Pseudomonas infections.
Veterinary medicine
Fractures
Chinchillas live active lives and can recover well from minor physical injury. Fractures may be problematic, because chinchillas sit on their hind legs and eat with their front paws, so many types of injuries will disturb their natural eating behavior.
Convulsions
Chinchilla breeders sometimes report seeing their animals have convulsions. Typically this happens only irregularly and then only for a few seconds, and not more than a few minutes at the most. Convulsions are a symptom that can have many causes, including a brain problem such as hemorrhaging, a vitamin or dietary element deficiency in the diet, or some kind of nervous system injury. If convulsions are observed after chinchillas mate then it is likely related to a circulatory problem. Some chinchillas who are kept in groups have stress convulsions during feeding if they see other chinchillas getting food first. Vitamin B, cardiac medication, or a calcium injection may be used to prevent convulsions.
Infectious diseases
Listeriosis is not a typical chinchilla disease, but rats in group housing conditions it can spread as a digestive tract disease in a community. Pasteurella can be contracted from food and then transmitted among a group of chinchillas. Symptoms include apathy, digestive disorder, and fever. Pseudomonas aeruginosa infections are widely distributed in nature and can affect chinchillas like many other animals. They can cause wide deaths in populations of chinchillas and spontaneous abortion in pregnant chinchillas.
Respiratory tract infections can be caused by many pathogens, but, regardless of cause, usually result in difficult breathing and a nasal discharge. Young chinchilla are more likely to be affected and these infections are unlikely to result in an epidemic, even if transmissible.
Gastrointestinal disorders are observed as either constipation or diarrhea. These are almost always the result of a problem with the diet, but if the diet is optimal, they could be the symptom of an infectious disease. Constipation in chinchillas is difficult to observe in groups because it may not be obvious that an animal is not contributing to the population's waste. If it is identified, mild treatments include feeding paraffin to soften the feces.
Mental health
Chinchillas are easily distressed, and when they are unhappy, they may exhibit physical symptoms. A common indicator of stress in pet chinchillas is fur-chewing (or fur barbering), an excessive grooming behavior that results in uneven patches of fur; chinchillas may chew their own fur or that of their cagemates. Fur-chewing can sometimes be alleviated through changes in living environment, but is regarded by some experts to be passed genetically from parents to offspring. Usually, fur-chewing itself is a benign symptom that does not cause physiological distress.
Sick chinchillas may stop eating if they are stressed, which can make them even weaker. Chinchillas that live in communities are especially sensitive in their breeding seasons of February to March and August to September. Chinchillas are social animals and are likely to be upset to have their breeding mate changed in breeding season. They are known to be disturbed by a change of diet in these times.
Pharmaceutical treatment
Chinchillas may be treated with chloramphenicol, neomycin, or spectinomycin for digestive problems. Sulfonamides dissolved in drinking water may be used. Colistin can be an effective antibiotic.
| Biology and health sciences | Rodents | Animals |
198319 | https://en.wikipedia.org/wiki/Hamiltonian%20mechanics | Hamiltonian mechanics | In physics, Hamiltonian mechanics is a reformulation of Lagrangian mechanics that emerged in 1833. Introduced by Sir William Rowan Hamilton, Hamiltonian mechanics replaces (generalized) velocities used in Lagrangian mechanics with (generalized) momenta. Both theories provide interpretations of classical mechanics and describe the same physical phenomena.
Hamiltonian mechanics has a close relationship with geometry (notably, symplectic geometry and Poisson structures) and serves as a link between classical and quantum mechanics.
Overview
Phase space coordinates (p, q) and Hamiltonian H
Let be a mechanical system with configuration space and smooth Lagrangian Select a standard coordinate system on The quantities are called momenta. (Also generalized momenta, conjugate momenta, and canonical momenta). For a time instant the Legendre transformation of is defined as the map which is assumed to have a smooth inverse For a system with degrees of freedom, the Lagrangian mechanics defines the energy function
The Legendre transform of turns into a function known as the . The Hamiltonian satisfies
which implies that
where the velocities are found from the (-dimensional) equation which, by assumption, is uniquely solvable for . The (-dimensional) pair is called phase space coordinates. (Also canonical coordinates).
From Euler–Lagrange equation to Hamilton's equations
In phase space coordinates , the (-dimensional) Euler–Lagrange equation
becomes Hamilton's equations in dimensions
From stationary action principle to Hamilton's equations
Let be the set of smooth paths for which and The action functional is defined via
where , and (see above). A path is a stationary point of (and hence is an equation of motion) if and only if the path in phase space coordinates obeys the Hamilton's equations.
Basic physical interpretation
A simple interpretation of Hamiltonian mechanics comes from its application on a one-dimensional system consisting of one nonrelativistic particle of mass . The value of the Hamiltonian is the total energy of the system, in this case the sum of kinetic and potential energy, traditionally denoted and , respectively. Here is the momentum and is the space coordinate. Then
is a function of alone, while is a function of alone (i.e., and are scleronomic).
In this example, the time derivative of is the velocity, and so the first Hamilton equation means that the particle's velocity equals the derivative of its kinetic energy with respect to its momentum. The time derivative of the momentum equals the Newtonian force, and so the second Hamilton equation means that the force equals the negative gradient of potential energy.
Example
A spherical pendulum consists of a mass m moving without friction on the surface of a sphere. The only forces acting on the mass are the reaction from the sphere and gravity. Spherical coordinates are used to describe the position of the mass in terms of , where is fixed, .
The Lagrangian for this system is
Thus the Hamiltonian is
where
and
In terms of coordinates and momenta, the Hamiltonian reads
Hamilton's equations give the time evolution of coordinates and conjugate momenta in four first-order differential equations,
Momentum , which corresponds to the vertical component of angular momentum , is a constant of motion. That is a consequence of the rotational symmetry of the system around the vertical axis. Being absent from the Hamiltonian, azimuth is a cyclic coordinate, which implies conservation of its conjugate momentum.
Deriving Hamilton's equations
Hamilton's equations can be derived by a calculation with the Lagrangian , generalized positions , and generalized velocities , where . Here we work off-shell, meaning , , are independent coordinates in phase space, not constrained to follow any equations of motion (in particular, is not a derivative of ). The total differential of the Lagrangian is:
The generalized momentum coordinates were defined as , so we may rewrite the equation as:
After rearranging, one obtains:
The term in parentheses on the left-hand side is just the Hamiltonian defined previously, therefore:
One may also calculate the total differential of the Hamiltonian with respect to coordinates , , instead of , , , yielding:
One may now equate these two expressions for , one in terms of , the other in terms of :
Since these calculations are off-shell, one can equate the respective coefficients of , , on the two sides:
On-shell, one substitutes parametric functions which define a trajectory in phase space with velocities , obeying Lagrange's equations:
Rearranging and writing in terms of the on-shell gives:
Thus Lagrange's equations are equivalent to Hamilton's equations:
In the case of time-independent and , i.e. , Hamilton's equations consist of first-order differential equations, while Lagrange's equations consist of second-order equations. Hamilton's equations usually do not reduce the difficulty of finding explicit solutions, but important theoretical results can be derived from them, because coordinates and momenta are independent variables with nearly symmetric roles.
Hamilton's equations have another advantage over Lagrange's equations: if a system has a symmetry, so that some coordinate does not occur in the Hamiltonian (i.e. a cyclic coordinate), the corresponding momentum coordinate is conserved along each trajectory, and that coordinate can be reduced to a constant in the other equations of the set. This effectively reduces the problem from coordinates to coordinates: this is the basis of symplectic reduction in geometry. In the Lagrangian framework, the conservation of momentum also follows immediately, however all the generalized velocities still occur in the Lagrangian, and a system of equations in coordinates still has to be solved.
The Lagrangian and Hamiltonian approaches provide the groundwork for deeper results in classical mechanics, and suggest analogous formulations in quantum mechanics: the path integral formulation and the Schrödinger equation.
Properties of the Hamiltonian
The value of the Hamiltonian is the total energy of the system if and only if the energy function has the same property. (See definition of ).
when , form a solution of Hamilton's equations. Indeed, and everything but the final term cancels out.
does not change under point transformations, i.e. smooth changes of space coordinates. (Follows from the invariance of the energy function under point transformations. The invariance of can be established directly).
(See ).
. (Compare Hamilton's and Euler-Lagrange equations or see ).
if and only if .A coordinate for which the last equation holds is called cyclic (or ignorable). Every cyclic coordinate reduces the number of degrees of freedom by , causes the corresponding momentum to be conserved, and makes Hamilton's equations easier to solve.
Hamiltonian as the total system energy
In its application to a given system, the Hamiltonian is often taken to be
where is the kinetic energy and is the potential energy. Using this relation can be simpler than first calculating the Lagrangian, and then deriving the Hamiltonian from the Lagrangian. However, the relation is not true for all systems.
The relation holds true for nonrelativistic systems when all of the following conditions are satisfied
where is time, is the number of degrees of freedom of the system, and each is an arbitrary scalar function of .
In words, this means that the relation holds true if does not contain time as an explicit variable (it is scleronomic), does not contain generalised velocity as an explicit variable, and each term of is quadratic in generalised velocity.
Proof
Preliminary to this proof, it is important to address an ambiguity in the related mathematical notation. While a change of variables can be used to equate
,
it is important to note that
.
In this case, the right hand side always evaluates to 0. To perform a change of variables inside of a partial derivative, the multivariable chain rule should be used. Hence, to avoid ambiguity, the function arguments of any term inside of a partial derivative should be stated.
Additionally, this proof uses the notation to imply that .
Application to systems of point masses
For a system of point masses, the requirement for to be quadratic in generalised velocity is always satisfied for the case where , which is a requirement for anyway.
Conservation of energy
If the conditions for are satisfied, then conservation of the Hamiltonian implies conservation of energy. This requires the additional condition that does not contain time as an explicit variable.
With respect to the extended Euler-Lagrange formulation (See ), the Rayleigh dissipation function represents energy dissipation by nature. Therefore, energy is not conserved when . This is similar to the velocity dependent potential.
In summary, the requirements for to be satisfied for a nonrelativistic system are
is a homogeneous quadratic function in
Hamiltonian of a charged particle in an electromagnetic field
A sufficient illustration of Hamiltonian mechanics is given by the Hamiltonian of a charged particle in an electromagnetic field. In Cartesian coordinates the Lagrangian of a non-relativistic classical particle in an electromagnetic field is (in SI Units):
where is the electric charge of the particle, is the electric scalar potential, and the are the components of the magnetic vector potential that may all explicitly depend on and .
This Lagrangian, combined with Euler–Lagrange equation, produces the Lorentz force law
and is called minimal coupling.
The canonical momenta are given by:
The Hamiltonian, as the Legendre transformation of the Lagrangian, is therefore:
This equation is used frequently in quantum mechanics.
Under gauge transformation:
where is any scalar function of space and time. The aforementioned Lagrangian, the canonical momenta, and the Hamiltonian transform like:
which still produces the same Hamilton's equation:
In quantum mechanics, the wave function will also undergo a local U(1) group transformation during the Gauge Transformation, which implies that all physical results must be invariant under local U(1) transformations.
Relativistic charged particle in an electromagnetic field
The relativistic Lagrangian for a particle (rest mass and charge ) is given by:
Thus the particle's canonical momentum is
that is, the sum of the kinetic momentum and the potential momentum.
Solving for the velocity, we get
So the Hamiltonian is
This results in the force equation (equivalent to the Euler–Lagrange equation)
from which one can derive
The above derivation makes use of the vector calculus identity:
An equivalent expression for the Hamiltonian as function of the relativistic (kinetic) momentum, , is
This has the advantage that kinetic momentum can be measured experimentally whereas canonical momentum cannot. Notice that the Hamiltonian (total energy) can be viewed as the sum of the relativistic energy (kinetic+rest), , plus the potential energy, .
From symplectic geometry to Hamilton's equations
Geometry of Hamiltonian systems
The Hamiltonian can induce a symplectic structure on a smooth even-dimensional manifold in several equivalent ways, the best known being the following:
As a closed nondegenerate symplectic 2-form ω. According to the Darboux's theorem, in a small neighbourhood around any point on there exist suitable local coordinates (canonical or symplectic coordinates) in which the symplectic form becomes:
The form induces a natural isomorphism of the tangent space with the cotangent space: . This is done by mapping a vector to the 1-form , where for all . Due to the bilinearity and non-degeneracy of , and the fact that , the mapping is indeed a linear isomorphism. This isomorphism is natural in that it does not change with change of coordinates on Repeating over all , we end up with an isomorphism between the infinite-dimensional space of smooth vector fields and that of smooth 1-forms. For every and ,
(In algebraic terms, one would say that the -modules and are isomorphic). If , then, for every fixed , , and . is known as a Hamiltonian vector field. The respective differential equation on
is called . Here and is the (time-dependent) value of the vector field at .
A Hamiltonian system may be understood as a fiber bundle over time , with the fiber being the position space at time . The Lagrangian is thus a function on the jet bundle over ; taking the fiberwise Legendre transform of the Lagrangian produces a function on the dual bundle over time whose fiber at is the cotangent space , which comes equipped with a natural symplectic form, and this latter function is the Hamiltonian. The correspondence between Lagrangian and Hamiltonian mechanics is achieved with the tautological one-form.
Any smooth real-valued function on a symplectic manifold can be used to define a Hamiltonian system. The function is known as "the Hamiltonian" or "the energy function." The symplectic manifold is then called the phase space. The Hamiltonian induces a special vector field on the symplectic manifold, known as the Hamiltonian vector field.
The Hamiltonian vector field induces a Hamiltonian flow on the manifold. This is a one-parameter family of transformations of the manifold (the parameter of the curves is commonly called "the time"); in other words, an isotopy of symplectomorphisms, starting with the identity. By Liouville's theorem, each symplectomorphism preserves the volume form on the phase space. The collection of symplectomorphisms induced by the Hamiltonian flow is commonly called "the Hamiltonian mechanics" of the Hamiltonian system.
The symplectic structure induces a Poisson bracket. The Poisson bracket gives the space of functions on the manifold the structure of a Lie algebra.
If and are smooth functions on then the smooth function is properly defined; it is called a Poisson bracket of functions and and is denoted . The Poisson bracket has the following properties:
bilinearity
antisymmetry
Leibniz rule:
Jacobi identity:
non-degeneracy: if the point on is not critical for then a smooth function exists such that .
Given a function
if there is a probability distribution , then (since the phase space velocity has zero divergence and probability is conserved) its convective derivative can be shown to be zero and so
This is called Liouville's theorem. Every smooth function over the symplectic manifold generates a one-parameter family of symplectomorphisms and if , then is conserved and the symplectomorphisms are symmetry transformations.
A Hamiltonian may have multiple conserved quantities . If the symplectic manifold has dimension and there are functionally independent conserved quantities which are in involution (i.e., ), then the Hamiltonian is Liouville integrable. The Liouville–Arnold theorem says that, locally, any Liouville integrable Hamiltonian can be transformed via a symplectomorphism into a new Hamiltonian with the conserved quantities as coordinates; the new coordinates are called action–angle coordinates. The transformed Hamiltonian depends only on the , and hence the equations of motion have the simple form
for some function . There is an entire field focusing on small deviations from integrable systems governed by the KAM theorem.
The integrability of Hamiltonian vector fields is an open question. In general, Hamiltonian systems are chaotic; concepts of measure, completeness, integrability and stability are poorly defined.
Riemannian manifolds
An important special case consists of those Hamiltonians that are quadratic forms, that is, Hamiltonians that can be written as
where is a smoothly varying inner product on the fibers , the cotangent space to the point in the configuration space, sometimes called a cometric. This Hamiltonian consists entirely of the kinetic term.
If one considers a Riemannian manifold or a pseudo-Riemannian manifold, the Riemannian metric induces a linear isomorphism between the tangent and cotangent bundles. (See Musical isomorphism). Using this isomorphism, one can define a cometric. (In coordinates, the matrix defining the cometric is the inverse of the matrix defining the metric.) The solutions to the Hamilton–Jacobi equations for this Hamiltonian are then the same as the geodesics on the manifold. In particular, the Hamiltonian flow in this case is the same thing as the geodesic flow. The existence of such solutions, and the completeness of the set of solutions, are discussed in detail in the article on geodesics. | Physical sciences | Classical mechanics | null |
198476 | https://en.wikipedia.org/wiki/Triangulum%20Galaxy | Triangulum Galaxy | The Triangulum Galaxy is a spiral galaxy 2.73 million light-years (ly) from Earth in the constellation Triangulum. It is catalogued as Messier 33 or NGC 598. With the D25 isophotal diameter of , the Triangulum Galaxy is the third-largest member of the Local Group of galaxies, behind the Andromeda Galaxy and the Milky Way.
The galaxy is the second-smallest spiral galaxy in the Local Group after the Large Magellanic Cloud, which is a Magellanic-type spiral galaxy. It is believed to be a satellite of the Andromeda Galaxy or on its rebound into the latter due to their interactions, velocities, and proximity to one another in the night sky. It also has an H II nucleus.
Etymology
The galaxy gets its name from the constellation Triangulum, where it can be spotted.
It is sometimes informally referred to as the "Pinwheel Galaxy" by some astronomy references, in some computerized telescope software, and in some public outreach websites. However, the SIMBAD Astronomical Database, a professional database, collates formal designations for astronomical objects and indicates that Pinwheel Galaxy refers to Messier 101, which several amateur astronomy resources including public outreach websites identify by that name, and that is within the bounds of Ursa Major.
Visibility
Under exceptionally good viewing conditions with no light pollution, the Triangulum Galaxy can be seen by some people with the fully dark-adapted naked eye; to those viewers, it is the farthest permanent entity visible without magnification, being about half again as distant as Messier 31, the Andromeda Galaxy. It is a diffuse, or extended, object rather than a starlike point, even without magnification, because of its physical extent.
Its observability without optical aid ranges from being relatively easily seen by people using direct vision in deep rural locations under a dark, clear, transparent sky, to requiring use of averted vision by observers in locations beyond the suburbs in shallow rural areas under good viewing conditions. It is one of the reference objects of the Bortle Dark-Sky Scale.
Crumey has shown that although the total apparent V-magnitude of M33 is 5.72, it has an effective visual magnitude of approximately 6.6, meaning that a precondition for visibility is that the observer can see stars at least as faint as that latter figure. This is fainter than many people are able to see, even at a very dark site.
Observation history
The Triangulum Galaxy was probably discovered by the Italian astronomer Giovanni Battista Hodierna before 1654. In his work De systemate orbis cometici; deque admirandis coeli caracteribus ("About the systematics of the cometary orbit, and about the admirable objects of the sky"), he listed it as a cloud-like nebulosity or obscuration and gave the cryptic description, "near the Triangle hinc inde". This is in reference to the constellation Triangulum as a pair of triangles. The magnitude of the object matches M33, so it is most likely a reference to the Triangulum Galaxy.
The galaxy was independently discovered by Charles Messier on the night of August 25–26, 1764. It was published in his Catalog of Nebulae and Star Clusters (1771) as object number 33; hence the name M33. When William Herschel compiled his extensive catalog of nebulae, he was careful not to include most of the objects identified by Messier. However, M33 was an exception, and he cataloged this object on September 11, 1784, as H V-17.
Herschel also cataloged the Triangulum Galaxy's brightest and largest H II region (diffuse emission nebula containing ionized hydrogen) as H III.150 separately from the galaxy itself; the nebula eventually obtained NGC number 604. As seen from Earth, NGC 604 is located northeast of the galaxy's central core. It is one of the largest H II regions known, with a diameter of nearly 1500 light-years and a spectrum similar to that of the Orion Nebula. Herschel also noted three other smaller H II regions (NGC 588, 592, and 595).
It was among the first "spiral nebulae" identified as such by Lord Rosse in 1850. In 1922–23, John Charles Duncan and Max Wolf discovered variable stars in the nebulae. Edwin Hubble showed in 1926 that 35 of these stars were classical Cepheids, thereby allowing him to estimate their distances. The results were consistent with the concept of spiral nebulae being independent galactic systems of gas and dust, rather than just nebulae in the Milky Way.
Properties
The Triangulum Galaxy is the third largest member of the Local Group of galaxies. It has a diameter measured through the D25 standard - the isophote where the surface brightness of the galaxy reaches 25 mag/arcsec2, to be about , making it roughly 70% the size of the Milky Way. It may be a gravitationally bound companion of the Andromeda Galaxy. Triangulum may be home to 40 billion stars, compared to 400 billion for the Milky Way and 1 trillion for Andromeda.
The disk of Triangulum has an estimated mass of solar masses, while the gas component is about solar masses. Thus, the combined mass of all baryonic matter in the galaxy may be 1010 solar masses. The contribution of the dark matter component out to a radius of is equivalent to about solar masses.
Location – distance – motion
Estimates of the distance from the Milky Way to the Triangulum Galaxy range from (or 2.38 to 3.07 Mly), with most estimates since the year 2000 lying in the middle portion of this range, making it slightly more distant than the Andromeda Galaxy (at 2,540,000 light-years). At least three techniques have been used to measure distances to M 33. Using the Cepheid variable method, an estimate of was achieved in 2004. In the same year, the tip of the red-giant branch (TRGB) method was used to derive a distance estimate of . The Triangulum Galaxy is around 750,000 light years from the Andromeda Galaxy.
In 2006, a group of astronomers announced the discovery of an eclipsing binary star in the Triangulum Galaxy. By studying the eclipses of the stars, astronomers were able to measure their sizes. Knowing the sizes and temperatures of the stars, they were able to measure the absolute magnitude of the stars. When the visual and absolute magnitudes are known, the distance to the star can be measured. The stars lie at the distance of . The average of 102 distance estimates published since 1987 gives a distance modulus of 24.69, or .883 Mpc (2,878,000 light-years).
The Triangulum Galaxy is a source of H2O maser emission. In 2005, using observations of two water masers on opposite sides of Triangulum via the VLBA, researchers were for the first time able to estimate the angular rotation and proper motion of Triangulum. A velocity of relative to the Milky Way was computed, which means Triangulum is moving towards Andromeda Galaxy and suggesting it may be a satellite of the larger galaxy (depending on their relative distances and margins of error).
In 2004, evidence was announced of a clumpy stream of hydrogen gas linking the Andromeda Galaxy with Triangulum, suggesting that the two may have tidally interacted in the past. This discovery was confirmed in 2011. A distance of less than 300 kiloparsecs between the two supports this hypothesis.
The Pisces Dwarf (LGS 3), one of the small Local Group member galaxies, is located from the Sun. It is 20° from the Andromeda Galaxy and 11° from Triangulum. As LGS 3 lies at a distance of from both galaxies, it could be a satellite galaxy of either Andromeda or Triangulum. LGS 3 has a core radius of and solar masses.
Pisces VII/Triangulum (Tri) III may be another satellite of Triangulum.
Structure
In the French astronomer Gérard de Vaucouleurs' revised Hubble Sandage (VRHS) system of galaxy morphological classification, the Triangulum Galaxy is classified as type SA(s)cd. The S prefix indicates that it is a disk-shaped galaxy with prominent arms of gas and dust that spiral out from the nucleus—what is commonly known as a spiral galaxy. The A is assigned when the galactic nucleus lacks a bar-shaped structure, in contrast to SB class barred spiral galaxies. American astronomer Allan Sandage's "(s)" notation is used when the spiral arms emerge directly from the nucleus or central bar, rather than from an inner ring as with an (r)-type galaxy. Finally, the cd suffix represents a stage along the spiral sequence that describes the openness of the arms. A rating of cd indicates relatively loosely wound arms.
This galaxy has an inclination of 54° to the line of sight from Earth, allowing the structure to be examined without significant obstruction by gas and dust. The disk of the Triangulum Galaxy appears warped out to a radius of about 8 kpc. There may be a halo surrounding the galaxy, but there is no bulge at the nucleus. This is an isolated galaxy and there are no indications of recent mergers or interactions with other galaxies, and it lacks the dwarf spheroidals or tidal tails associated with the Milky Way.
Triangulum is classified as unbarred, but an analysis of the galaxy's shape shows what may be a weak bar-like structure about the galactic nucleus. The radial extent of this structure is about 0.8 kpc.
The nucleus of this galaxy is an H II region, and it contains an ultraluminous X-ray source with an emission of , which is the most luminous source of X-rays in the Local Group of galaxies. This source is modulated by 20% over a 106-day cycle. However, the nucleus does not appear to contain a supermassive black hole, as a best-fit value of zero mass and an upper limit of is placed on the mass of a central black hole based on models and the Hubble Space Telescope (HST) data. This is significantly lower than the mass expected from the velocity dispersion of the nucleus and far below any mass predicted from the disk kinematics. This may suggest that supermassive black holes are associated only with galaxy bulges instead of with their disks. Assuming that the upper limit of the central black hole is correct, it would be rather an intermediate-mass black hole.
The inner part of the galaxy has two luminous spiral arms, along with multiple spurs that connect the inner to the outer spiral features. The main arms are designated IN (north) and IS (south).
Star formation
In the central 4′ region of this galaxy, atomic gas is being efficiently converted to molecular gas, resulting in a strong spectral emission of CO. This effect occurs as giant molecular clouds condense out of the surrounding interstellar medium. A similar process is taking place outside the central 4′, but at a less efficient pace. About 10% of the gas content in this galaxy is in the molecular form.
Star formation is taking place at a rate that is strongly correlated with local gas density, and the rate per unit area is higher than in the neighboring Andromeda Galaxy. (The rate of star formation is about 3.4 solar masses Gyr−1 pc−2 in the Triangulum Galaxy, compared to 0.74 in Andromeda.) The total integrated rate of star formation in the Triangulum Galaxy is about . It is uncertain whether this net rate is currently decreasing or remaining constant.
Based on analysis of the chemical composition of this galaxy, it appears to be divided into two distinct components with differing histories. The inner disk within a radius of has a typical composition gradient that decreases linearly from the core. Beyond this radius, out to about , the gradient is much flatter. This suggests a different star formation history between the inner disk and the outer disk and halo, and may be explained by a scenario of "inside-out" galaxy formation. This occurs when gas is accumulated at large radii later in a galaxy's life space, while the gas at the core becomes exhausted. The result is a decrease in the average age of stars with increasing radius from the galaxy core.
Discrete features
Using infrared observations from the Spitzer Space Telescope, a total of 515 discrete candidate sources of 24 μm emission within the Triangulum Galaxy have been catalogued as of 2007. The brightest sources lie within the central region of the galaxy and along the spiral arms.
Many of the emission sources are associated with H II regions of star formation. The four brightest HII regions are designated NGC 588, NGC 592, NGC 595, and NGC 604. These regions are associated with molecular clouds containing solar masses. The brightest of these regions, NGC 604, may have undergone a discrete outburst of star formation about three million years ago. This nebula is the second most luminous HII region within the Local Group of galaxies, at times the luminosity of the Sun. Other prominent HII regions in Triangulum include IC 132, IC 133, and IK 53.
The northern main spiral arm contains four large HII regions, while the southern arm has greater concentrations of young, hot stars. The estimated rate of supernova explosions in the Triangulum Galaxy is 0.06 Type Ia and 0.62 Type Ib/Type II per century. This is equivalent to a supernova explosion every 147 years, on average. As of 2008, a total of 100 supernova remnants have been identified in the Triangulum Galaxy, the majority of which lie in the southern half of the spiral galaxy. Similar asymmetries exist for H I and H II regions, plus highly luminous concentrations of massive, O type stars. The center of the distribution of these features is offset about two arc minutes to the southwest. M33 being a local galaxy, the Central Bureau for Astronomical Telegrams (CBAT) tracks novae in it along with M31 and M81.
About 54 globular clusters have been identified in this galaxy, but the actual number may be 122 or more. The confirmed clusters may be several billion years younger than globular clusters in the Milky Way, and cluster formation appears to have increased during the past 100 million years. This increase is correlated with an inflow of gas into the center of the galaxy. The ultraviolet emission of massive stars in this galaxy matches the level of similar stars in the Large Magellanic Cloud.
In 2007, a black hole about 15.7 times the mass of the Sun was detected in this galaxy using data from the Chandra X-ray Observatory. The black hole, named M33 X-7, orbits a companion star which it eclipses every 3.5 days. It is the largest stellar mass black hole known.
Unlike the Milky Way and Andromeda galaxies, the Triangulum Galaxy does not appear to have a supermassive black hole at its center. This may be because the mass of a galaxy's central supermassive black hole correlates with the size of the galaxy's central bulge, and unlike the Milky Way and Andromeda, the Triangulum Galaxy is a pure disk galaxy with no bulge.
Relationship with the Andromeda Galaxy
As mentioned above, M33 is linked to M31 by several streams of neutral hydrogen and stars, which suggests that a past interaction between these two galaxies took place from 2 to 8 billion years ago, and a more violent encounter will occur 2.5 billion years in the future.
The fate of M33 was uncertain in 2009 beyond seeming to be linked to its larger neighbor M31. Suggested scenarios include being torn apart and absorbed by the greater companion, fueling the latter with hydrogen to form new stars; eventually exhausting all of its gas, and thus the ability to form new stars; or participating in the collision between the Milky Way and M31, likely ending up orbiting the merger product and fusing with it much later. Two other possibilities are a collision with the Milky Way before the Andromeda Galaxy arrives or an ejection out of the Local Group. Astrometric data from Gaia appears in 2019 to rule out the possibility that M33 and M31 are in orbit. If correct, M33 is on its first infall proper into the Andromeda Galaxy (M31).
| Physical sciences | Notable galaxies | null |
198491 | https://en.wikipedia.org/wiki/Leishmaniasis | Leishmaniasis | Leishmaniasis is a wide array of clinical manifestations caused by protozoal parasites of the Trypanosomatida genus Leishmania. It is generally spread through the bite of phlebotomine sandflies, Phlebotomus and Lutzomyia, and occurs most frequently in the tropics and sub-tropics of Africa, Asia, the Americas, and southern Europe. The disease can present in three main ways: cutaneous, mucocutaneous, or visceral. The cutaneous form presents with skin ulcers, while the mucocutaneous form presents with ulcers of the skin, mouth, and nose. The visceral form starts with skin ulcers and later presents with fever, low red blood cell count, and enlarged spleen and liver.
Infections in humans are caused by more than 20 species of Leishmania. Risk factors include poverty, malnutrition, deforestation, and urbanization. All three types can be diagnosed by seeing the parasites under microscopy. Additionally, visceral disease can be diagnosed by blood tests.
Leishmaniasis can be partly prevented by sleeping under nets treated with insecticide. Other measures include spraying insecticides to kill sandflies and treating people with the disease early to prevent further spread. The treatment needed is determined by where the disease is acquired, the species of Leishmania, and the type of infection. Some possible medications used for visceral disease include liposomal amphotericin B, a combination of pentavalent antimonials and paromomycin, and miltefosine. For cutaneous disease, paromomycin, fluconazole, or pentamidine may be effective.
About 4 to 12 million people are currently infected in some 98 countries. About 2 million new cases and between 20 and 50 thousand deaths occur each year. About 200 million people in Asia, Africa, South and Central America, and southern Europe live in areas where the disease is common. The World Health Organization has obtained discounts on some medications to treat the disease. It is classified as a neglected tropical disease. The disease may occur in a number of other animals, including dogs and rodents.
Signs and symptoms
The symptoms of leishmaniasis are skin sores which erupt weeks to months after the person is bitten by infected sandflies.
Leishmaniasis may be divided into the following types:
Cutaneous leishmaniasis is the most common form, which causes an open sore at each bite site, which heals in a few months to a year and a half, leaving an unpleasant-looking scar.
Diffuse cutaneous leishmaniasis produces widespread skin lesions which resemble leprosy, and may not heal on their own.
Mucocutaneous leishmaniasis causes both skin and mucosal ulcers with damage primarily of the nose and mouth.
Visceral leishmaniasis or kala-azar ('black fever') is the most serious form and is generally fatal if untreated. Other consequences, which can occur a few months to years after infection, include fever, damage to the spleen and liver, and anemia.
Leishmaniasis is considered one of the classic causes of a markedly enlarged (and therefore palpable) spleen; the organ, which is not normally felt during the examination of the abdomen, may even become larger than the liver in severe cases.
Cause
Leishmaniasis is transmitted by the bite of infected female phlebotomine sandflies which can transmit the protozoa Leishmania. The sandflies inject the infective stage, metacyclic promastigotes, during blood meals. Metacyclic promastigotes in the puncture wound are phagocytized by macrophages, and transform into amastigotes. Amastigotes multiply in infected cells and affect different tissues, depending in part on the host, and in part on which Leishmania species is involved. These differing tissue specificities cause the differing clinical manifestations of the various forms of leishmaniasis. Sandflies become infected during blood meals on infected hosts when they ingest macrophages infected with amastigotes. In the sandfly's midgut, the parasites differentiate into promastigotes, which multiply, differentiate into metacyclic promastigotes, and migrate to the proboscis.
The genomes of three Leishmania species (L. major, L. infantum, and L. braziliensis) have been sequenced, and this has provided much information about the biology of the parasite. For example, in Leishmania, protein-coding genes are understood to be organized as large polycistronic units in a head-to-head or tail-to-tail manner; RNA polymerase II transcribes long polycistronic messages in the absence of defined RNA pol II promoters, and Leishmania has unique features concerning the regulation of gene expression in response to changes in the environment. The new knowledge from these studies may help identify new targets for urgently needed drugs and aid the development of vaccines.
Vector
Although most of the literature mentions only one genus transmitting Leishmania to humans (Lutzomyia) in the New World, a 2003 study by Galati suggested a new classification for New World sand flies, elevating several subgenera to the genus level. Elsewhere in the world, the genus Phlebotomus is considered the vector of leishmaniasis.
Possible non-human reservoirs
Some cases of infection of non-human animals of human-infecting species of Leishmania have been observed. In one study, L. major was identified in twelve out of ninety-one wild western lowland gorilla fecal samples and in a study of fifty-two captive non-human primates under zoo captivity in a leishmaniasis endemic area, eight (all three chimpanzees, three golden lion tamarins, a tufted capuchin, and an Angolan talapoin), were found to be infected with L. infantum and capable of infecting Lutzomyia longipalpis sand flies, although "parasite loads in infected sand flies observed in this study were considered low".
Organisms
Visceral disease is usually caused by Leishmania donovani, L. infantum, or L. chagasi, but occasionally these species may cause other forms of disease. The cutaneous form of the disease is caused by more than 15 species of Leishmania.
Risk factors
Risk factors include malnutrition, deforestation, lack of sanitation, suppressed immune system, and urbanization.
Diagnosis
Leishmaniasis is diagnosed in the hematology laboratory by direct visualization of the amastigotes (Leishman–Donovan bodies). Buffy-coat preparations of peripheral blood or aspirates from marrow, spleen, lymph nodes, or skin lesions should be spread on a slide to make a thin smear and stained with Leishman stain or Giemsa stain (pH 7.2) for 20 minutes. Amastigotes are seen within blood and spleen monocytes or, less commonly, in circulating neutrophils and in aspirated tissue macrophages. They are small, round bodies 2–4 μm in diameter with indistinct cytoplasm, a nucleus, and a small, rod-shaped kinetoplast. Occasionally, amastigotes may be seen lying free between cells. However, the retrieval of tissue samples is often painful for the patient and identification of the infected cells can be difficult. So, other indirect immunological methods of diagnosis are developed, including enzyme-linked immunosorbent assay, antigen-coated dipsticks, and direct agglutination test. Although these tests are readily available, they are not the standard diagnostic tests due to their insufficient sensitivity and specificity.
Several different polymerase chain reaction (PCR) tests are available for the detection of Leishmania DNA. With this assay, a specific and sensitive diagnostic procedure is finally possible. The most sensitive PCR tests use minicircle kinetoplast DNA found in the parasite. Kinetoplast DNA contains sequences for mitochondrial proteins in its maxicircles(~25–50 per parasite), and guide RNA in its minicircles(~10'000 per parasite) of the kinetoplast. With this specific method, one can still detect Leishmania even with a very low parasite load. When needing to diagnose a specific species of Leishmania, as opposed to only detection, other PCR methods have been superior.
Most forms of the disease are transmitted only from nonhuman animals, but some can be spread between humans. Infections in humans are caused by about 21 of 30 species that infect mammals; the different species look the same, but they can be differentiated by isoenzyme analysis, DNA sequence analysis, or monoclonal antibodies.
Prevention
Using insect repellent on exposed skin and under the ends of sleeves and pant legs. Follow the instructions on the label of the repellent. The most effective repellents generally are those that contain the chemical DEET (N,N-diethylmetatoluamide)
Leishmaniasis can be partly prevented by using nets treated with insecticide or insect repellent while sleeping. To provide good protection against sandflies, fine mesh sizes of 0.6 mm or less are required, but a mosquito net with 1.2mm mesh will provide a limited reduction in the number of sandfly bites. Finer mesh sizes have the downside of higher cost and reduced air circulation which can cause overheating. Many Phlebotomine sandfly attacks occur at sunset rather than at night, so it may also be useful to put nets over doors and windows or to use insect repellents.
Use of insecticide-impregnated dog collars and treatment or culling of infected dogs.
Spraying houses and animal shelters with insecticides.
Treatment
The treatment is determined by where the disease is acquired, the species of Leishmania, and the type of infection.
For visceral leishmaniasis in India, South America, and the Mediterranean, liposomal amphotericin B is the recommended treatment and is often used as a single dose. Rates of cure with a single dose of amphotericin have been reported as 95%. In India, almost all infections are resistant to pentavalent antimonials. In Africa, a combination of pentavalent antimonials and paromomycin is recommended. These, however, can have significant side effects. Miltefosine, an oral medication, is effective against both visceral and cutaneous leishmaniasis. Side effects are generally mild, though it can cause birth defects if taken within three months of getting pregnant. It does not appear to work for L. major or L. braziliensis. Trifluralin, a herbicide, is shown to be effective treatment as ointment, without hemolytic or cell-toxic side-effects.
The evidence around the treatment of cutaneous leishmaniasis is poor. Several topical treatments may be used for cutaneous leishmaniasis. Which treatments are effective depends on the strain, with topical paromomycin effective for L. major, L. tropica, L. mexicana, L. panamensis, and L. braziliensis. Pentamidine is effective for L. guyanensis. Oral fluconazole or itraconazole appears effective in L. major and L. tropica. There is limited evidence to support the use of heat therapy in cutaneous leishmaniasis as of 2015.
No studies have determined the effect of oral nutritional supplements on visceral leishmaniasis being treated with anti-leishmanial drug therapy.
Epidemiology
Out of 200 countries and territories reporting to WHO, 97 countries and territories are endemic for leishmaniasis. The settings in which leishmaniasis is found range from rainforests in Central and South America to deserts in western Asia and the Middle East. It affects as many as 12 million people worldwide, with 1.5–2.0 million new cases each year. The visceral form of leishmaniasis has an estimated incidence of 500,000 new
cases. In 2014, more than 90% of new cases reported to WHO occurred in six countries: Brazil, Ethiopia, India, Somalia, South Sudan and Sudan. it caused about 52,000 deaths, down from 87,000 in 1990. Different types of the disease occur in different regions of the world. Cutaneous disease is most common in Afghanistan, Algeria, Brazil, Colombia, and Iran, while mucocutaneous disease is most common in Bolivia, Brazil, and Peru, and visceral disease is most common in Bangladesh, Brazil, Ethiopia, India, and Sudan.
Leishmaniasis is found through much of the Americas from northern Argentina to South Texas, though not in Uruguay or Chile, and has recently been shown to be spreading to North Texas and Oklahoma, and further expansion to the north may be facilitated by climate change as more habitat becomes suitable for vector and reservoir species for leishmaniasis. Leishmaniasis is also known as papalomoyo, papa lo moyo, úlcera de los chicleros, and chiclera in Latin America. During 2004, an estimated 3,400 troops from the Colombian army, operating in the jungles near the south of the country (in particular around the Meta and Guaviare departments), were infected with leishmaniasis. Allegedly, a contributing factor was that many of the affected soldiers did not use the officially provided insect repellent because of its disturbing odor. Nearly 13,000 cases of the disease were recorded in all of Colombia throughout 2004, and about 360 new instances of the disease among soldiers had been reported in February 2005.
The disease is found across much of Asia and in the Middle East. Within Afghanistan, leishmaniasis occurs commonly in Kabul, partly due to bad sanitation and waste left uncollected in streets, allowing parasite-spreading sand flies an environment they find favorable. In Kabul, the number of people infected was estimated to be at least 200,000, and in three other towns (Herat, Kandahar, and Mazar-i-Sharif) about 70,000 more occurred, according to WHO figures from 2002. Kabul is estimated as the largest center of cutaneous leishmaniasis in the world, with around 67,500 cases as of 2004. Africa, in particular, the East and North, is also home to cases of leishmaniasis. Leishmaniasis is considered endemic also in some parts of southern parts of western Europe and has spread towards the north in recent years. For example, an outbreak of cutaneous and visceral leishmaniasis was reported from Madrid, Spain, between 2010 and 2012.
Leishmaniasis is mostly a disease of the developing world and is rarely known in the developed world outside a small number of cases, mostly in instances where troops are stationed away from their home countries. Leishmaniasis has been reported by U.S. troops stationed in Saudi Arabia and Iraq since the Gulf War of 1990, including visceral leishmaniasis.
In September 2005, the disease was contracted by at least four Dutch marines who were stationed in Mazar-i-Sharif, Afghanistan, and subsequently repatriated for treatment.
History
Descriptions of conspicuous lesions similar to cutaneous leishmaniasis appear on tablets from King Ashurbanipal from the seventh century BCE, some of which may have derived from even earlier texts from 1500 to 2500 BCE. Persian physicians, including Avicenna in the 10th century CE, gave detailed descriptions of what was called balkh sore. In 1756, Alexander Russell, after examining a Turkish patient, gave one of the most detailed clinical descriptions of the disease. Physicians in the Indian subcontinent would describe it as kala-azar (pronounced kālā āzār, the Urdu, Hindi, and Hindustani phrase for "black fever", kālā meaning black and āzār meaning fever or disease). In the Americas, evidence of the cutaneous form of the disease in Ecuador and Peru appears in pre-Inca pottery depicting skin lesions and deformed faces dating back to the first century CE. Some 15th- and 16th-century texts from the Inca period and from Spanish colonials mention "valley sickness", "Andean sickness", or "white leprosy", which are likely to be the cutaneous form.
It remains unclear who first discovered the organism. David Douglas Cunningham, Surgeon Major of the British Indian army, may have seen it in 1885 without being able to relate it to the disease. Peter Borovsky, a Russian military surgeon working in Tashkent, conducted research into the etiology of "oriental sore", locally known as sart sore, and in 1898 published the first accurate description of the causative agent, correctly described the parasite's relation to host tissues and correctly referred it to the protozoa. However, because his results were published in Russian in a journal with low circulation, his results were not internationally acknowledged during his lifetime. In 1901, William Boog Leishman identified certain organisms in smears taken from the spleen of a patient who had died from "dum-dum fever" (Dum Dum is an area close to Calcutta) and proposed them to be trypanosomes, found for the first time in India. A few months later, Captain Charles Donovan (1863–1951) confirmed the finding of what became known as Leishman-Donovan bodies in smears taken from people in Madras in southern India. But it was Ronald Ross who proposed that Leishman-Donovan bodies were the intracellular stages of a new parasite, which he named Leishmania donovani. The link with the disease kala-azar was first suggested by Charles Donovan, and was conclusively demonstrated by Charles Bentley's discovery of L. donovani in patients with kala-azar. Transmission by the sandfly was hypothesized by Lionel Napier and Ernest Struthers at the School of Tropical Medicine at Calcutta and later proven by his colleagues. The disease became a major problem for Allied troops fighting in Sicily during the Second World War; research by Leonard Goodwin then showed pentostam was an effective treatment.
Society and culture
The Institute for OneWorld Health has reintroduced the drug paromomycin for the treatment of leishmaniasis, results which led to its approval as an orphan drug. The Drugs for Neglected Diseases Initiative is also actively facilitating the search for novel therapeutics. A treatment with paromomycin will cost about US$10. The drug had originally been identified in the 1950s but had been abandoned because it would not be profitable, as the disease mostly affects poor people. The Indian government approved paromomycin for sale in August 2006.
By 2012 the World Health Organization had successfully negotiated with the manufacturers to achieve a reduced cost for liposomal amphotericin B, to US$18 a vial, but several vials are needed for treatment and it must be kept at a stable, cool temperature.
Research
As of 2017, no leishmaniasis vaccine for humans was available. Research to produce a human vaccine is ongoing.
Currently some effective leishmaniasis vaccines for dogs exist. There is also the consideration that public health practices can control or eliminate leishmaniasis without a vaccine. Pyrimidine–based drugs are being explored as anti-leishmanial compounds.
| Biology and health sciences | Protozoan infections | Health |
198507 | https://en.wikipedia.org/wiki/Scientific%20theory | Scientific theory | A scientific theory is an explanation of an aspect of the natural world and universe that can be or that has been repeatedly tested and has corroborating evidence in accordance with the scientific method, using accepted protocols of observation, measurement, and evaluation of results. Where possible, theories are tested under controlled conditions in an experiment. In circumstances not amenable to experimental testing, theories are evaluated through principles of abductive reasoning. Established scientific theories have withstood rigorous scrutiny and embody scientific knowledge.
A scientific theory differs from a scientific fact or scientific law in that a theory seeks to explain "why" or "how", whereas a fact is a simple, basic observation and a law is an empirical description of a relationship between facts and/or other laws. For example, Newton's Law of Gravity is a mathematical equation that can be used to predict the attraction between bodies, but it is not a theory to explain how gravity works. Stephen Jay Gould wrote that "...facts and theories are different things, not rungs in a hierarchy of increasing certainty. Facts are the world's data. Theories are structures of ideas that explain and interpret facts."
The meaning of the term scientific theory (often contracted to theory for brevity) as used in the disciplines of science is significantly different from the common vernacular usage of theory. In everyday speech, theory can imply an explanation that represents an unsubstantiated and speculative guess, whereas in a scientific context it most often refers to an explanation that has already been tested and is widely accepted as valid.
The strength of a scientific theory is related to the diversity of phenomena it can explain and its simplicity. As additional scientific evidence is gathered, a scientific theory may be modified and ultimately rejected if it cannot be made to fit the new findings; in such circumstances, a more accurate theory is then required. Some theories are so well-established that they are unlikely ever to be fundamentally changed (for example, scientific theories such as evolution, heliocentric theory, cell theory, theory of plate tectonics, germ theory of disease, etc.). In certain cases, a scientific theory or scientific law that fails to fit all data can still be useful (due to its simplicity) as an approximation under specific conditions. An example is Newton's laws of motion, which are a highly accurate approximation to special relativity at velocities that are small relative to the speed of light.
Scientific theories are testable and make verifiable predictions. They describe the causes of a particular natural phenomenon and are used to explain and predict aspects of the physical universe or specific areas of inquiry (for example, electricity, chemistry, and astronomy). As with other forms of scientific knowledge, scientific theories are both deductive and inductive, aiming for predictive and explanatory power. Scientists use theories to further scientific knowledge, as well as to facilitate advances in technology or medicine. Scientific hypotheses can never be "proven" because scientists are not able to fully confirm that their hypothesis is true. Instead, scientists say that the study "supports" or is consistent with their hypothesis.
Types
Albert Einstein described two different types of scientific theories: "Constructive theories" and "principle theories". Constructive theories are constructive models for phenomena: for example, kinetic theory. Principle theories are empirical generalisations, one such example being Newton's laws of motion.
Characteristics
Essential criteria
For any theory to be accepted within most academia there is usually one simple criterion. The essential criterion is that the theory must be observable and repeatable. The aforementioned criterion is essential to prevent fraud and perpetuate science itself.
The defining characteristic of all scientific knowledge, including theories, is the ability to make falsifiable or testable predictions. The relevance and specificity of those predictions determine how potentially useful the theory is. A would-be theory that makes no observable predictions is not a scientific theory at all. Predictions not sufficiently specific to be tested are similarly not useful. In both cases, the term "theory" is not applicable.
A body of descriptions of knowledge can be called a theory if it fulfills the following criteria:
It makes falsifiable predictions with consistent accuracy across a broad area of scientific inquiry (such as mechanics).
It is well-supported by many independent strands of evidence, rather than a single foundation.
It is consistent with preexisting experimental results and at least as accurate in its predictions as are any preexisting theories.
These qualities are certainly true of such established theories as special and general relativity, quantum mechanics, plate tectonics, the modern evolutionary synthesis, etc.
Other criteria
In addition, most scientists prefer to work with a theory that meets the following qualities:
It can be subjected to minor adaptations to account for new data that do not fit it perfectly, as they are discovered, thus increasing its predictive capability over time.
It is among the most parsimonious explanations, economical in the use of proposed entities or explanatory steps as per Occam's razor. This is because for each accepted explanation of a phenomenon, there may be an extremely large, perhaps even incomprehensible, number of possible and more complex alternatives, because one can always burden failing explanations with ad hoc hypotheses to prevent them from being falsified; therefore, simpler theories are preferable to more complex ones because they are more testable.
Definitions from scientific organizations
The United States National Academy of Sciences defines scientific theories as follows:
The formal scientific definition of theory is quite different from the everyday meaning of the word. It refers to a comprehensive explanation of some aspect of nature that is supported by a vast body of evidence. Many scientific theories are so well established that no new evidence is likely to alter them substantially. For example, no new evidence will demonstrate that the Earth does not orbit around the Sun (heliocentric theory), or that living things are not made of cells (cell theory), that matter is not composed of atoms, or that the surface of the Earth is not divided into solid plates that have moved over geological timescales (the theory of plate tectonics)...One of the most useful properties of scientific theories is that they can be used to make predictions about natural events or phenomena that have not yet been observed.
From the American Association for the Advancement of Science:
A scientific theory is a well-substantiated explanation of some aspect of the natural world, based on a body of facts that have been repeatedly confirmed through observation and experiment. Such fact-supported theories are not "guesses" but reliable accounts of the real world. The theory of biological evolution is more than "just a theory". It is as factual an explanation of the universe as the atomic theory of matter or the germ theory of disease. Our understanding of gravity is still a work in progress. But the phenomenon of gravity, like evolution, is an accepted fact.
Note that the term theory would not be appropriate for describing untested but intricate hypotheses or even scientific models.
Formation
The scientific method involves the proposal and testing of hypotheses, by deriving predictions from the hypotheses about the results of future experiments, then performing those experiments to see whether the predictions are valid. This provides evidence either for or against the hypothesis. When enough experimental results have been gathered in a particular area of inquiry, scientists may propose an explanatory framework that accounts for as many of these as possible. This explanation is also tested, and if it fulfills the necessary criteria (see above), then the explanation becomes a theory. This can take many years, as it can be difficult or complicated to gather sufficient evidence.
Once all of the criteria have been met, it will be widely accepted by scientists (see scientific consensus) as the best available explanation of at least some phenomena. It will have made predictions of phenomena that previous theories could not explain or could not predict accurately, and it will have many repeated bouts of testing. The strength of the evidence is evaluated by the scientific community, and the most important experiments will have been replicated by multiple independent groups.
Theories do not have to be perfectly accurate to be scientifically useful. For example, the predictions made by classical mechanics are known to be inaccurate in the relativistic realm, but they are almost exactly correct at the comparatively low velocities of common human experience. In chemistry, there are many acid-base theories providing highly divergent explanations of the underlying nature of acidic and basic compounds, but they are very useful for predicting their chemical behavior. Like all knowledge in science, no theory can ever be completely certain, since it is possible that future experiments might conflict with the theory's predictions. However, theories supported by the scientific consensus have the highest level of certainty of any scientific knowledge; for example, that all objects are subject to gravity or that life on Earth evolved from a common ancestor.
Acceptance of a theory does not require that all of its major predictions be tested, if it is already supported by sufficient evidence. For example, certain tests may be unfeasible or technically difficult. As a result, theories may make predictions that have not yet been confirmed or proven incorrect; in this case, the predicted results may be described informally with the term "theoretical". These predictions can be tested at a later time, and if they are incorrect, this may lead to the revision or rejection of the theory.As Feynman puts it:It doesn't matter how beautiful your theory is, it doesn't matter how smart you are. If it doesn't agree with experiment, it's wrong.
Modification and improvement
If experimental results contrary to a theory's predictions are observed, scientists first evaluate whether the experimental design was sound, and if so they confirm the results by independent replication. A search for potential improvements to the theory then begins. Solutions may require minor or major changes to the theory, or none at all if a satisfactory explanation is found within the theory's existing framework. Over time, as successive modifications build on top of each other, theories consistently improve and greater predictive accuracy is achieved. Since each new version of a theory (or a completely new theory) must have more predictive and explanatory power than the last, scientific knowledge consistently becomes more accurate over time.
If modifications to the theory or other explanations seem to be insufficient to account for the new results, then a new theory may be required. Since scientific knowledge is usually durable, this occurs much less commonly than modification. Furthermore, until such a theory is proposed and accepted, the previous theory will be retained. This is because it is still the best available explanation for many other phenomena, as verified by its predictive power in other contexts. For example, it has been known since 1859 that the observed perihelion precession of Mercury violates Newtonian mechanics, but the theory remained the best explanation available until relativity was supported by sufficient evidence. Also, while new theories may be proposed by a single person or by many, the cycle of modifications eventually incorporates contributions from many different scientists.
After the changes, the accepted theory will explain more phenomena and have greater predictive power (if it did not, the changes would not be adopted); this new explanation will then be open to further replacement or modification. If a theory does not require modification despite repeated tests, this implies that the theory is very accurate. This also means that accepted theories continue to accumulate evidence over time, and the length of time that a theory (or any of its principles) remains accepted often indicates the strength of its supporting evidence.
Unification
In some cases, two or more theories may be replaced by a single theory that explains the previous theories as approximations or special cases, analogous to the way a theory is a unifying explanation for many confirmed hypotheses; this is referred to as unification of theories. For example, electricity and magnetism are now known to be two aspects of the same phenomenon, referred to as electromagnetism.
When the predictions of different theories appear to contradict each other, this is also resolved by either further evidence or unification. For example, physical theories in the 19th century implied that the Sun could not have been burning long enough to allow certain geological changes as well as the evolution of life. This was resolved by the discovery of nuclear fusion, the main energy source of the Sun. Contradictions can also be explained as the result of theories approximating more fundamental (non-contradictory) phenomena. For example, atomic theory is an approximation of quantum mechanics. Current theories describe three separate fundamental phenomena of which all other theories are approximations; the potential unification of these is sometimes called the Theory of Everything.
Example: Relativity
In 1905, Albert Einstein published the principle of special relativity, which soon became a theory. Special relativity predicted the alignment of the Newtonian principle of Galilean invariance, also termed Galilean relativity, with the electromagnetic field. By omitting from special relativity the luminiferous aether, Einstein stated that time dilation and length contraction measured in an object in relative motion is inertial—that is, the object exhibits constant velocity, which is speed with direction, when measured by its observer. He thereby duplicated the Lorentz transformation and the Lorentz contraction that had been hypothesized to resolve experimental riddles and inserted into electrodynamic theory as dynamical consequences of the aether's properties. An elegant theory, special relativity yielded its own consequences, such as the equivalence of mass and energy transforming into one another and the resolution of the paradox that an excitation of the electromagnetic field could be viewed in one reference frame as electricity, but in another as magnetism.
Einstein sought to generalize the invariance principle to all reference frames, whether inertial or accelerating. Rejecting Newtonian gravitation—a central force acting instantly at a distance—Einstein presumed a gravitational field. In 1907, Einstein's equivalence principle implied that a free fall within a uniform gravitational field is equivalent to inertial motion. By extending special relativity's effects into three dimensions, general relativity extended length contraction into space contraction, conceiving of 4D space-time as the gravitational field that alters geometrically and sets all local objects' pathways. Even massless energy exerts gravitational motion on local objects by "curving" the geometrical "surface" of 4D space-time. Yet unless the energy is vast, its relativistic effects of contracting space and slowing time are negligible when merely predicting motion. Although general relativity is embraced as the more explanatory theory via scientific realism, Newton's theory remains successful as merely a predictive theory via instrumentalism. To calculate trajectories, engineers and NASA still uses Newton's equations, which are simpler to operate.
Theories and laws
Both scientific laws and scientific theories are produced from the scientific method through the formation and testing of hypotheses, and can predict the behavior of the natural world. Both are also typically well-supported by observations and/or experimental evidence. However, scientific laws are descriptive accounts of how nature will behave under certain conditions. Scientific theories are broader in scope, and give overarching explanations of how nature works and why it exhibits certain characteristics. Theories are supported by evidence from many different sources, and may contain one or several laws.
A common misconception is that scientific theories are rudimentary ideas that will eventually graduate into scientific laws when enough data and evidence have been accumulated. A theory does not change into a scientific law with the accumulation of new or better evidence. A theory will always remain a theory; a law will always remain a law. Both theories and laws could potentially be falsified by countervailing evidence.
Theories and laws are also distinct from hypotheses. Unlike hypotheses, theories and laws may be simply referred to as scientific fact.
However, in science, theories are different from facts even when they are well supported. For example, evolution is both a theory and a fact.
About theories
Theories as axioms
The logical positivists thought of scientific theories as statements in a formal language. First-order logic is an example of a formal language. The logical positivists envisaged a similar scientific language. In addition to scientific theories, the language also included observation sentences ("the sun rises in the east"), definitions, and mathematical statements. The phenomena explained by the theories, if they could not be directly observed by the senses (for example, atoms and radio waves), were treated as theoretical concepts. In this view, theories function as axioms: predicted observations are derived from the theories much like theorems are derived in Euclidean geometry. However, the predictions are then tested against reality to verify the predictions, and the "axioms" can be revised as a direct result.
The phrase "the received view of theories" is used to describe this approach. Terms commonly associated with it are "linguistic" (because theories are components of a language) and "syntactic" (because a language has rules about how symbols can be strung together). Problems in defining this kind of language precisely, e.g., are objects seen in microscopes observed or are they theoretical objects, led to the effective demise of logical positivism in the 1970s.
Theories as models
The semantic view of theories, which identifies scientific theories with models rather than propositions, has replaced the received view as the dominant position in theory formulation in the philosophy of science. A model is a logical framework intended to represent reality (a "model of reality"), similar to the way that a map is a graphical model that represents the territory of a city or country.
In this approach, theories are a specific category of models that fulfill the necessary criteria (see above). One can use language to describe a model; however, the theory is the model (or a collection of similar models), and not the description of the model. A model of the solar system, for example, might consist of abstract objects that represent the sun and the planets. These objects have associated properties, e.g., positions, velocities, and masses. The model parameters, e.g., Newton's Law of Gravitation, determine how the positions and velocities change with time. This model can then be tested to see whether it accurately predicts future observations; astronomers can verify that the positions of the model's objects over time match the actual positions of the planets. For most planets, the Newtonian model's predictions are accurate; for Mercury, it is slightly inaccurate and the model of general relativity must be used instead.
The word "semantic" refers to the way that a model represents the real world. The representation (literally, "re-presentation") describes particular aspects of a phenomenon or the manner of interaction among a set of phenomena. For instance, a scale model of a house or of a solar system is clearly not an actual house or an actual solar system; the aspects of an actual house or an actual solar system represented in a scale model are, only in certain limited ways, representative of the actual entity. A scale model of a house is not a house; but to someone who wants to learn about houses, analogous to a scientist who wants to understand reality, a sufficiently detailed scale model may suffice.
Differences between theory and model
Several commentators have stated that the distinguishing characteristic of theories is that they are explanatory as well as descriptive, while models are only descriptive (although still predictive in a more limited sense). Philosopher Stephen Pepper also distinguished between theories and models, and said in 1948 that general models and theories are predicated on a "root" metaphor that constrains how scientists theorize and model a phenomenon and thus arrive at testable hypotheses.
Engineering practice makes a distinction between "mathematical models" and "physical models"; the cost of fabricating a physical model can be minimized by first creating a mathematical model using a computer software package, such as a computer aided design tool. The component parts are each themselves modelled, and the fabrication tolerances are specified. An exploded view drawing is used to lay out the fabrication sequence. Simulation packages for displaying each of the subassemblies allow the parts to be rotated, magnified, in realistic detail. Software packages for creating the bill of materials for construction allows subcontractors to specialize in assembly processes, which spreads the cost of manufacturing machinery among multiple customers. See: Computer-aided engineering, Computer-aided manufacturing, and 3D printing
Assumptions in formulating theories
An assumption (or axiom) is a statement that is accepted without evidence. For example, assumptions can be used as premises in a logical argument. Isaac Asimov described assumptions as follows:
...it is incorrect to speak of an assumption as either true or false, since there is no way of proving it to be either (If there were, it would no longer be an assumption). It is better to consider assumptions as either useful or useless, depending on whether deductions made from them corresponded to reality...Since we must start somewhere, we must have assumptions, but at least let us have as few assumptions as possible.
Certain assumptions are necessary for all empirical claims (e.g. the assumption that reality exists). However, theories do not generally make assumptions in the conventional sense (statements accepted without evidence). While assumptions are often incorporated during the formation of new theories, these are either supported by evidence (such as from previously existing theories) or the evidence is produced in the course of validating the theory. This may be as simple as observing that the theory makes accurate predictions, which is evidence that any assumptions made at the outset are correct or approximately correct under the conditions tested.
Conventional assumptions, without evidence, may be used if the theory is only intended to apply when the assumption is valid (or approximately valid). For example, the special theory of relativity assumes an inertial frame of reference. The theory makes accurate predictions when the assumption is valid, and does not make accurate predictions when the assumption is not valid. Such assumptions are often the point with which older theories are succeeded by new ones (the general theory of relativity works in non-inertial reference frames as well).
The term "assumption" is actually broader than its standard use, etymologically speaking. The Oxford English Dictionary (OED) and online Wiktionary indicate its Latin source as assumere ("accept, to take to oneself, adopt, usurp"), which is a conjunction of ad- ("to, towards, at") and sumere (to take). The root survives, with shifted meanings, in the Italian assumere and Spanish sumir. The first sense of "assume" in the OED is "to take unto (oneself), receive, accept, adopt". The term was originally employed in religious contexts as in "to receive up into heaven", especially "the reception of the Virgin Mary into heaven, with body preserved from corruption", (1297 CE) but it was also simply used to refer to "receive into association" or "adopt into partnership". Moreover, other senses of assumere included (i) "investing oneself with (an attribute)", (ii) "to undertake" (especially in Law), (iii) "to take to oneself in appearance only, to pretend to possess", and (iv) "to suppose a thing to be" (all senses from OED entry on "assume"; the OED entry for "assumption" is almost perfectly symmetrical in senses). Thus, "assumption" connotes other associations than the contemporary standard sense of "that which is assumed or taken for granted; a supposition, postulate" (only the 11th of 12 senses of "assumption", and the 10th of 11 senses of "assume").
Descriptions
From philosophers of science
Karl Popper described the characteristics of a scientific theory as follows:
It is easy to obtain confirmations, or verifications, for nearly every theory—if we look for confirmations.
Confirmations should count only if they are the result of risky predictions; that is to say, if, unenlightened by the theory in question, we should have expected an event which was incompatible with the theory—an event which would have refuted the theory.
Every "good" scientific theory is a prohibition: it forbids certain things to happen. The more a theory forbids, the better it is.
A theory which is not refutable by any conceivable event is non-scientific. Irrefutability is not a virtue of a theory (as people often think) but a vice.
Every genuine test of a theory is an attempt to falsify it, or to refute it. Testability is falsifiability; but there are degrees of testability: some theories are more testable, more exposed to refutation, than others; they take, as it were, greater risks.
Confirming evidence should not count except when it is the result of a genuine test of the theory; and this means that it can be presented as a serious but unsuccessful attempt to falsify the theory. (I now speak in such cases of "corroborating evidence".)
Some genuinely testable theories, when found to be false, might still be upheld by their admirers—for example by introducing post hoc (after the fact) some auxiliary hypothesis or assumption, or by reinterpreting the theory post hoc in such a way that it escapes refutation. Such a procedure is always possible, but it rescues the theory from refutation only at the price of destroying, or at least lowering, its scientific status, by tampering with evidence. The temptation to tamper can be minimized by first taking the time to write down the testing protocol before embarking on the scientific work.
Popper summarized these statements by saying that the central criterion of the scientific status of a theory is its "falsifiability, or refutability, or testability". Echoing this, Stephen Hawking states, "A theory is a good theory if it satisfies two requirements: It must accurately describe a large class of observations on the basis of a model that contains only a few arbitrary elements, and it must make definite predictions about the results of future observations." He also discusses the "unprovable but falsifiable" nature of theories, which is a necessary consequence of inductive logic, and that "you can disprove a theory by finding even a single observation that disagrees with the predictions of the theory".
Several philosophers and historians of science have, however, argued that Popper's definition of theory as a set of falsifiable statements is wrong because, as Philip Kitcher has pointed out, if one took a strictly Popperian view of "theory", observations of Uranus when first discovered in 1781 would have "falsified" Newton's celestial mechanics. Rather, people suggested that another planet influenced Uranus' orbit—and this prediction was indeed eventually confirmed.
Kitcher agrees with Popper that "There is surely something right in the idea that a science can succeed only if it can fail." He also says that scientific theories include statements that cannot be falsified, and that good theories must also be creative. He insists we view scientific theories as an "elaborate collection of statements", some of which are not falsifiable, while others—those he calls "auxiliary hypotheses", are.
According to Kitcher, good scientific theories must have three features:
Unity: "A science should be unified.... Good theories consist of just one problem-solving strategy, or a small family of problem-solving strategies, that can be applied to a wide range of problems."
Fecundity: "A great scientific theory, like Newton's, opens up new areas of research.... Because a theory presents a new way of looking at the world, it can lead us to ask new questions, and so to embark on new and fruitful lines of inquiry.... Typically, a flourishing science is incomplete. At any time, it raises more questions than it can currently answer. But incompleteness is not vice. On the contrary, incompleteness is the mother of fecundity.... A good theory should be productive; it should raise new questions and presume those questions can be answered without giving up its problem-solving strategies."
Auxiliary hypotheses that are independently testable: "An auxiliary hypothesis ought to be testable independently of the particular problem it is introduced to solve, independently of the theory it is designed to save." (For example, the evidence for the existence of Neptune is independent of the anomalies in Uranus's orbit.)
Like other definitions of theories, including Popper's, Kitcher makes it clear that a theory must include statements that have observational consequences. But, like the observation of irregularities in the orbit of Uranus, falsification is only one possible consequence of observation. The production of new hypotheses is another possible and equally important result.
Analogies and metaphors
The concept of a scientific theory has also been described using analogies and metaphors. For example, the logical empiricist Carl Gustav Hempel likened the structure of a scientific theory to a "complex spatial network:"
Its terms are represented by the knots, while the threads connecting the latter correspond, in part, to the definitions and, in part, to the fundamental and derivative hypotheses included in the theory. The whole system floats, as it were, above the plane of observation and is anchored to it by the rules of interpretation. These might be viewed as strings which are not part of the network but link certain points of the latter with specific places in the plane of observation. By virtue of these interpretive connections, the network can function as a scientific theory: From certain observational data, we may ascend, via an interpretive string, to some point in the theoretical network, thence proceed, via definitions and hypotheses, to other points, from which another interpretive string permits a descent to the plane of observation.
Michael Polanyi made an analogy between a theory and a map:
A theory is something other than myself. It may be set out on paper as a system of rules, and it is the more truly a theory the more completely it can be put down in such terms. Mathematical theory reaches the highest perfection in this respect. But even a geographical map fully embodies in itself a set of strict rules for finding one's way through a region of otherwise uncharted experience. Indeed, all theory may be regarded as a kind of map extended over space and time.
A scientific theory can also be thought of as a book that captures the fundamental information about the world, a book that must be researched, written, and shared. In 1623, Galileo Galilei wrote:
Philosophy [i.e. physics] is written in this grand book—I mean the universe—which stands continually open to our gaze, but it cannot be understood unless one first learns to comprehend the language and interpret the characters in which it is written. It is written in the language of mathematics, and its characters are triangles, circles, and other geometrical figures, without which it is humanly impossible to understand a single word of it; without these, one is wandering around in a dark labyrinth.
The book metaphor could also be applied in the following passage, by the contemporary philosopher of science Ian Hacking:
I myself prefer an Argentine fantasy. God did not write a Book of Nature of the sort that the old Europeans imagined. He wrote a Borgesian library, each book of which is as brief as possible, yet each book of which is inconsistent with every other. No book is redundant. For every book there is some humanly accessible bit of Nature such that that book, and no other, makes possible the comprehension, prediction and influencing of what is going on...Leibniz said that God chose a world which maximized the variety of phenomena while choosing the simplest laws. Exactly so: but the best way to maximize phenomena and have simplest laws is to have the laws inconsistent with each other, each applying to this or that but none applying to all.
In physics
In physics, the term theory is generally used for a mathematical framework—derived from a small set of basic postulates (usually symmetries—like equality of locations in space or in time, or identity of electrons, etc.)—that is capable of producing experimental predictions for a given category of physical systems. A good example is classical electromagnetism, which encompasses results derived from gauge symmetry (sometimes called gauge invariance) in a form of a few equations called Maxwell's equations. The specific mathematical aspects of classical electromagnetic theory are termed "laws of electromagnetism", reflecting the level of consistent and reproducible evidence that supports them. Within electromagnetic theory generally, there are numerous hypotheses about how electromagnetism applies to specific situations. Many of these hypotheses are already considered to be adequately tested, with new ones always in the making and perhaps untested. An example of the latter might be the radiation reaction force. As of 2009, its effects on the periodic motion of charges are detectable in synchrotrons, but only as averaged effects over time. Some researchers are now considering experiments that could observe these effects at the instantaneous level (i.e. not averaged over time).
Examples
Note that many fields of inquiry do not have specific named theories, e.g. developmental biology. Scientific knowledge outside a named theory can still have a high level of certainty, depending on the amount of evidence supporting it. Also note that since theories draw evidence from many fields, the categorization is not absolute.
Biology: cell theory, theory of evolution (modern evolutionary synthesis), abiogenesis, germ theory, particulate inheritance theory, dual inheritance theory, Young–Helmholtz theory, opponent process, cohesion-tension theory
Chemistry: collision theory, kinetic theory of gases, Lewis theory, molecular theory, molecular orbital theory, transition state theory, valence bond theory
Physics: atomic theory, Big Bang theory, Dynamo theory, perturbation theory, theory of relativity (successor to classical mechanics), quantum field theory
Earth science: Climate change theory (from climatology), plate tectonics theory (from geology), theories of the origin of the Moon, theories for the Moon illusion
Astronomy: Self-gravitating system, Stellar evolution, solar nebular model, stellar nucleosynthesis
Explanatory notes
| Physical sciences | Basics | null |
198570 | https://en.wikipedia.org/wiki/Oil%20well | Oil well | An oil well is a drillhole boring in Earth that is designed to bring petroleum oil hydrocarbons to the surface. Usually some natural gas is released as associated petroleum gas along with the oil. A well that is designed to produce only gas may be termed a gas well. Wells are created by drilling down into an oil or gas reserve and if necessary equipped with extraction devices such as pumpjacks. Creating the wells can be an expensive process, costing at least hundreds of thousands of dollars, and costing much more when in difficult-to-access locations, e.g., offshore. The process of modern drilling for wells first started in the 19th century but was made more efficient with advances to oil drilling rigs and technology during the 20th century.
Wells are frequently sold or exchanged between different oil and gas companies as an asset – in large part because during falls in the price of oil and gas, a well may be unproductive, but if prices rise, even low-production wells may be economically valuable. Moreover, new methods, such as hydraulic fracturing (a process of injecting gas or liquid to force more oil or natural gas production) have made some wells viable. However, peak oil and climate policy surrounding fossil fuels have made fewer of these wells and costly techniques viable.
However, a large number of neglected or poorly maintained wellheads is a large environmental issue: they may leak methane or other toxic substances into local air, water and soil systems. This pollution often becomes worse when wells are abandoned or orphaned – i.e., where wells no longer economically viable are no longer maintained by their (former) owners. A 2020 estimate by Reuters suggested that there were at least 29 million abandoned wells internationally, creating a significant source of greenhouse gas emissions worsening climate change.
History
The earliest known oil wells were drilled in China in 347 CE. These wells had depths of up to about and were drilled using bits attached to bamboo poles. The oil was burned to evaporate brine producing salt. By the 10th century, extensive bamboo pipelines connected oil wells with salt springs. The ancient records of China and Japan are said to contain many allusions to the use of natural gas for lighting and heating. Petroleum was known as burning water in Japan in the 7th century.
According to Kasem Ajram, petroleum was distilled by the Persian alchemist Muhammad ibn Zakarīya Rāzi (Rhazes) in the 9th century, producing chemicals such as kerosene in the alembic (al-ambiq), and which was mainly used for kerosene lamps. Arab and Persian chemists also distilled crude oil in order to produce flammable products for military purposes. Through Islamic Spain, distillation became available in Western Europe by the 12th century.
Some sources claim that from the 9th century, oil fields were exploited in the area around modern Baku, Azerbaijan, to produce naphtha for the petroleum industry. These places were described by Marco Polo in the 13th century, who described the output of those oil wells as hundreds of shiploads. When Marco Polo in 1264 visited Baku, on the shores of the Caspian Sea, he saw oil being collected from seeps. He wrote that "on the confines toward Geirgine there is a fountain from which oil springs in great abundance, in as much as a hundred shiploads might be taken from it at one time."
In 1846, Baku (settlement Bibi-Heybat) the first ever well was drilled with percussion tools to a depth of for oil exploration. In 1846–1848, the first modern oil wells were drilled on the Absheron Peninsula north-east of Baku, by Russian engineer Vasily Semyonov applying the ideas of Nikolay Voskoboynikov.
Ignacy Łukasiewicz, a Polish pharmacist and petroleum industry pioneer drilled one of the world's first modern oil wells in 1854 in Polish village Bóbrka, Krosno County who in 1856 built one of the world's first oil refineries.
In North America, the first commercial oil well entered operation in Oil Springs, Ontario in 1858, while the first offshore oil well was drilled in 1896 in the Summerland Oil Field on the California Coast.
The earliest oil wells in modern times were drilled percussively, by repeatedly raising and dropping a bit on the bottom of a cable into the borehole. In the 20th century, cable tools were largely replaced with rotary drilling, which could drill boreholes to much greater depths and in less time. The record-depth Kola Borehole used a mud motor while drilling to achieve a depth of over .
Until the 1970s, most oil wells were essentially vertical, although lithological variations cause most wells to deviate at least slightly from true vertical (see deviation survey). However, modern directional drilling technologies allow for highly deviated wells that can, given sufficient depth and with the proper tools, actually become horizontal. This is of great value as the reservoir rocks that contain hydrocarbons are usually horizontal or nearly horizontal; a horizontal wellbore placed in a production zone has more surface area in the production zone than a vertical well, resulting in a higher production rate. The use of deviated and horizontal drilling has also made it possible to reach reservoirs several kilometers or miles away from the drilling location (extended reach drilling), allowing for the production of hydrocarbons located below locations that are difficult to place a drilling rig on, environmentally sensitive, or populated.
Life of a well
Planning
For a production well, the target is picked to optimize production from the well and manage reservoir drainage.
For an exploration or appraisal well, the target is chosen to confirm the existence of a viable hydrocarbon reservoir or to learn its extent.
For an injection well, the target is selected to locate the point of injection in a permeable zone that may support disposing of water or gas and/or pushing hydrocarbons into nearby production wells.
The target (the endpoint of the well) will be matched with a surface location (the starting point of the well), and a trajectory between the two will be designed. There are many considerations to take into account when designing the trajectory such as the clearance from any nearby wells (anti-collision) or future wellpaths.
Before a well is drilled, a geologic target is identified by a geologist or geophysicist to meet the objectives of the well.
When the well path is identified, a team of geoscientists and engineers will develop a set of presumed characteristics of the subsurface path that will be drilled through to reach the target. These properties may include lithology pore pressure, fracture gradient, wellbore stability, porosity and permeability. These assumptions are used by a well engineering team designing the casing and completion programs for the well. Also considered in the detailed planning are selection of the drill bits, Bottom hole assembly, and the drilling fluid Step-by-step procedures are written to provide guidelines for executing the well in a safe and cost-efficient manner.
With the interplay with many of the elements in a well's design, trajectories and designs often go through several iterations before the plan is finalized.
Drilling
The well is created by drilling a hole 12 cm to 1 meter (5 in to 40 in) in diameter into the earth with a drilling rig that rotates a drill string with a bit attached. At depths during the process, sections of steel pipe (casing), slightly smaller in diameter than the borehole at that point, are placed in the hole. Cement slurry will be pumped down the inside to rise in the annulus between the borehole and the outside of the casing. The casing provides structural integrity to that portion of the newly drilled wellbore, in addition to isolating potentially dangerous high pressure zones from lower-pressure ones, and from the surface.
With these zones safely isolated and the formation protected by the casing, the well can be drilled deeper (into potentially higher-pressure or more-unstable formations) with a smaller bit, and then cased with a smaller size pipe. Modern wells generally have two to as many as five sets of subsequently smaller hole sizes, each cemented with casing.
To drill the well
The rotating drill bit, aided by the weight of the drill string above it, cuts into the rock. There are different types of drill bits; some cause the rock to disintegrate by compressive failure, while others shear slices off the rock as the bit turns.
Drilling fluid, a.k.a. "mud", is pumped down the inside of the drill pipe and exits at the drill bit. The principal components of drilling fluid are usually water and clay, but it also typically contains a complex mixture of fluids, solids and chemicals that must be carefully tailored to provide the correct physical and chemical characteristics required to safely drill the well. Particular functions of the drilling mud include cooling the bit, lifting rock cuttings to the surface, preventing destabilisation (spalling) of the rock in the wellbore, and overcoming the pressure of fluids inside the rock so that these fluids do not enter the wellbore. Some oil wells are drilled with air or foam as the drilling fluid.
The generated rock "cuttings" are swept up by the drilling fluid as it circulates back to the surface inside the casing and outside of the drill pipe. The fluid then goes through "shakers" that screen the cuttings out of the fluid, which is returned to the pit for reuse. Watching for abnormalities in the returning cuttings and monitoring pit volume or rate of returning fluid are imperative to catch "kicks" early. A "kick" is when the formation pressure at the depth of the bit is greater than the hydrostatic head of the mud above, which if not controlled temporarily by closing the blowout preventers followed by increasing the density of the drilling fluid would allow formation fluids to enter the annulus uncontrollably.
The drill string to which the bit is attached is gradually lengthened as the well gets deeper by screwing in additional 9 m (30 ft) sections or "joints" of pipe under the kelly or top drive at the surface. This process is called "making a connection". The operation called "tripping" is when pulling the bit out of the hole to replace the bit (tripping out), and running back in with a new bit (tripping in). Joints are usually combined for more efficient tripping by creating stands of multiple joints. A conventional triple, for example, has three joints at a time racked vertically in the derrick. Some modern rigs, called "super singles", trip pipe one at a time, laying it out on racks as they go.
This process is all facilitated by a drilling rig, which contains all necessary equipment to circulate the drilling fluid, hoist and rotate the pipe, remove cuttings from the drilling fluid, and generate on-site power for these operations.
Completion
After drilling and casing the well, it must be 'completed'. Completion is the process in which the well is prepared to produce oil or gas.
In a cased-hole completion, small perforations are made in the portion of the casing across the production zone, to provide a path for the oil to flow from the surrounding rock into the production tubing. In open hole completion, often a 'sand screen' or 'gravel pack' is installed in the last-drilled but uncased reservoir section. These maintain structural integrity of the wellbore in the absence of casing, while still allowing flow from the reservoir into the borehole. Screens also control the migration of formation sands into production tubulars, which can lead to washouts and other problems, particularly from unconsolidated sand formations.
After a flow path is made, acids and fracturing fluids may be pumped into the well to fracture, clean, or otherwise prepare and stimulate the reservoir rock to allow optimal production of hydrocarbons into the wellbore. Usually the area above the producing section of the well is packed off inside the casing, and connected to the surface via a smaller diameter pipe called tubing. This arrangement provides a redundant barrier to leaks of hydrocarbons as well as allowing damaged sections to be replaced. Also, the smaller cross-sectional area of the tubing gives reservoir fluids an increased velocity to minimize liquid fallback that would create additional back pressure, and shields the casing from corrosive well fluids.
In many wells, the natural pressure of the subsurface reservoir is high enough for the oil or gas to flow to the surface. However, this is not always the case, especially in depleted fields where the pressures have been lowered by other producing wells, or in low-permeability oil reservoirs. Installing a smaller diameter tubing may be enough to help the production, but artificial lift methods may also be needed. Common solutions include surface pump jacks, downhole hydraulic pumps or gas lift assistance. Many new systems in recent years have been introduced for well completion. Multiple packer systems with frac ports or port collars in an all-in-one system have cut completion costs and improved production, especially in the case of horizontal wells. These new systems allow casing to run into the lateral zone equipped with proper packer/frac-port placement for optimal hydrocarbon recovery.
Production
The production stage is the most important stage of a well's life: when the oil and gas are produced. By this time, the oil rigs and workover rigs used to drill and complete the well will have moved off the wellbore, and the top is usually outfitted with a collection of valves called a Christmas tree or production tree. These valves regulate pressures, control flows, and allow access to the wellbore in case further completion work is needed. From the outlet valve of the production tree, the flow can be connected to a distribution network of pipelines and tanks to supply the product to refineries, natural gas compressor stations, or oil export terminals.
As long as the pressure in the reservoir remains high enough, the production tree is all that is required to produce the well. If the pressure depletes and it is considered economically viable, an artificial lift method mentioned in the completions section can be employed.
Workovers are often necessary in older wells, which may need smaller diameter tubing, scale or paraffin removal, acid matrix jobs, or completion in new zones of interest in a shallower reservoir. Such remedial work can be performed using workover rigs – also known as pulling units, completion rigs or "service rigs" – to pull and replace tubing, or by the use of well intervention techniques utilizing coiled tubing. Depending on the type of lift system and wellhead a rod rig or flushby can be used to change a pump without pulling the tubing.
Enhanced recovery methods such as water flooding, steam flooding, or CO2 flooding may be used to increase reservoir pressure and provide a "sweep" effect to push hydrocarbons out of the reservoir. Such methods require the use of injection wells (often chosen from old production wells in a carefully determined pattern), and are used when facing problems with reservoir pressure depletion or high oil viscosity, sometimes being employed early in a field's life. In certain cases – depending on the reservoir's geomechanics – reservoir engineers may determine that ultimate recoverable oil may be increased by applying a waterflooding strategy early in the field's development rather than later. Such enhanced recovery techniques are often called Secondary or "tertiary recovery".
Abandonment
Types of wells
By produced fluid
Wells that produce crude oil
Wells that produce crude oil and natural gas, or
Wells that only produce natural gas.
Natural gas, in a raw form known as associated petroleum gas, is almost always a by-product of producing oil. The short, light-gas carbon chains come out of solution when undergoing pressure reduction from the reservoir to the surface, similar to uncapping a bottle of soda where the carbon dioxide effervesces. If it escapes into the atmosphere intentionally it is known as vented gas, or if unintentionally as fugitive gas.
Unwanted natural gas can be a disposal problem at wells that are developed to produce oil. If there are no pipelines for natural gas near the wellhead it may be of no value to the oil well owner since it cannot reach the consumer markets. Such unwanted gas may then be burned off at the well site in a practice known as production flaring, but due to the energy resource waste and environmental damage concerns this practice is becoming less common.
Often, unwanted (or 'stranded' gas without a market) gas is returned back into the reservoir with an 'injection' well for storage or for re-pressurizing the producing formation. Another solution is to convert the natural gas to a liquid fuel. Gas to liquid (GTL) is a developing technology that converts stranded natural gas into synthetic gasoline, diesel or jet fuel through the Fischer–Tropsch process developed in World War II Germany. Like oil, such dense liquid fuels can be transported using conventional tankers for trucking to refineries or users. Proponents claim GTL fuels burn cleaner than comparable petroleum fuels. Most major international oil companies are in advanced development stages of GTL production, e.g. the Pearl GTL plant in Qatar, scheduled to come online in 2011. In locations such as the United States with a high natural gas demand, pipelines are usually favored to take the gas from the well site to the end consumer.
By location
Wells can be located:
Onshore, or
Offshore
Offshore wells can further be subdivided into
Wells with subsea wellheads, where the top of the well is sitting on the ocean floor under water, and often connected to a pipeline on the ocean floor.
Wells with 'dry' wellheads, where the top of the well is above the water on a platform or jacket, which also often contains processing equipment for the produced fluid.
While the location of the well will be a large factor in the type of equipment used to drill it, there is actually little downhole difference in the well itself. An offshore well targets a reservoir that happens to be underneath an ocean. Due to logistics and specialized equipment needed, drilling an offshore well is far more costly than a comparable onshore well. These wells dot the Southern and Central Great Plains, Southwestern United States, and are the most common wells in the Middle East.
By purpose
Another way to classify oil wells is by their purpose in contributing to the development of a resource. They can be characterized as:
wildcat wells that are drilled where little or no known geological information is available. The site may have been selected because of wells drilled some distance from the proposed location but to an underground structure that appeared similar to the proposed site. Individuals who drill wildcat wells are known as 'wildcatters'.
exploration wells are drilled purely for exploratory (information gathering) purposes in a new area. The site selection is usually based on seismic data, satellite surveys, etc. Details gathered in this well include the presence of hydrocarbon in the drilled location, the amount of fluid present and the depth at which oil or gas occurs.
appraisal wells may be needed to assess characteristics (such as flow rate, reservoir quantity) of a proven hydrocarbon accumulation. Such wells reduce uncertainty about the characteristics and properties of the hydrocarbon present in the field.
production wells are drilled primarily for producing oil or gas, once the producing structure and characteristics are determined.
development wells are wells drilled for the production of oil or gas already proven by appraisal drilling to be suitable for exploitation.
abandoned wells are wells permanently plugged in the drilling phase for technical reasons, or that had failed to locate commercially valuable hydrocarbons.
At a producing well site, active wells may be further categorized as:
oil producers producing predominantly liquid hydrocarbons, but most include some associated gas.
gas producers producing almost entirely gaseous hydrocarbons, consisting mostly of natural gas.
water injectors injecting water into the formation to maintain reservoir pressure, or simply to dispose of water produced with the hydrocarbons because even after treatment, it would be too oily and too saline to be considered clean for dumping overboard offshore, let alone into a fresh water resource in the case of onshore wells. Water injection into the producing zone frequently has a beneficial element of reservoir management; however, often produced water disposal is into shallower zones safely beneath any fresh water zones.
aquifer producers intentionally producing water for re-injection to manage pressure. If possible this water will come from the reservoir itself. Using aquifer produced water rather than water from other sources is to preclude chemical incompatibility that might lead to reservoir-plugging precipitates. These wells will generally be needed only if produced water from the oil or gas producers is insufficient for reservoir management purposes.
gas injectors injecting gas into the reservoir often as a means of disposal or sequestering for later production, but also to maintain reservoir pressure.
Lahee classification
New Field Wildcat (NFW) – far from other producing fields and on a structure that has not previously produced.
New Pool Wildcat (NPW) – new pools on already producing structure.
Deeper Pool Test (DPT) – on already producing structure and pool, but on a deeper pay zone.
Shallower Pool Test (SPT) – on already producing structure and pool, but on a shallower pay zone.
Outpost (OUT) – usually two or more locations from nearest productive area.
Development Well (DEV) – can be on the extension of a pay zone, or between existing wells (Infill).
Cost
The cost to drill a well depends mainly on the daily rate of the drilling rig, the extra services required to drill the well, the duration of the well program (including downtime and weather time), and the remoteness of the location (logistic supply costs).
The daily rates of offshore drilling rigs vary by their depth capability, and the market availability. Rig rates reported by industry web service show that the deepwater water floating drilling rigs are over twice the daily cost of the shallow water fleet, and rates for jack-up fleet can vary by factor of 3 depending upon capability.
With deepwater drilling rig rates in 2015 of around $520,000/day, and similar additional spread costs, a deepwater well of a duration of 100 days can cost around US$100 million.
With high-performance jackup rig rates in 2015 of around $177,000, and similar service costs, a high pressure, high-temperature well of duration 100 days can cost about US$30 million.
Onshore wells can be considerably cheaper, particularly if the field is at a shallow depth, where costs range from less than $4.9 million to $8.3 million, and the average completion costing $2.9 million to $5.6 million per well. Completion makes up a larger portion of onshore well costs than offshore wells, which generally have the added cost burden of a surface platform.
The total costs mentioned do not include the those associated with the risk of explosion and leakage of oil. Those costs include the cost of protecting against such disasters, the cost of the cleanup effort, and the hard-to-calculate cost of damage to the company's image.
Impacts on wildlife
The impacts of oil exploration and drilling are often irreversible, particularly for wildlife. Research indicates that caribou in Alaska show a marked avoidance of areas near oil wells and seismic lines due to disturbances. Drilling often destroys wildlife habitat, causing wildlife stress, and breaks up large areas into smaller isolated ones, changing the environment, and forcing animals to migrate elsewhere. It can also bring in new species that compete with or prey on existing animals. Even though the actual area taken up by oil and gas equipment might be small, negative effects can spread. Animals like mule deer and elk try to stay away from the noise and activity of drilling sites, sometimes moving miles away to find peace. This movement and avoidance can lead to less space for these animals affecting their numbers and health.
The Sage-grouse is another example of an animal that tries to avoid areas with drilling, which can lead to fewer of them surviving and reproducing. Different studies show that drilling in their habitats negatively impacts sage-grouse populations. In Wyoming, sage grouse studied between 1984 and 2008 show a roughly 2.5 percent annual population decline in males, correlating with the density of oil and gas wells. Factors such as sagebrush cover and precipitation seemed to have little effect on count changes. These results align with other studies highlighting the detrimental impact of oil and gas development on sage-grouse populations.
| Technology | Energy and fuel | null |
198584 | https://en.wikipedia.org/wiki/Laptop | Laptop | A laptop computer or notebook computer, also known as a laptop or notebook, is a small, portable personal computer (PC). Laptops typically have a clamshell form factor with a flat-panel screen on the inside of the upper lid and an alphanumeric keyboard and pointing device on the inside of the lower lid. Most of the computer's internal hardware is in the lower part, under the keyboard, although many modern laptops have a built-in webcam at the top of the screen, and some even feature a touchscreen display. In most cases, unlike tablet computers which run on mobile operating systems, laptops tend to run on desktop operating systems, which were originally developed for desktop computers.
Laptops are used in a variety of settings, such as at work (especially on business trips), in education, for playing games, web browsing, for personal multimedia, and for general home computer use. They can run on both AC power and rechargable battery packs and can be folded shut for convenient storage and transportation, making them suitable for mobile use. Laptops combine many of the input/output components and capabilities of a desktop computer into a single unit, including a display screen (usually in diagonal size), small speakers, a keyboard, and a pointing device (namely compact ones such as touchpads or pointing sticks). Hardware specifications may vary significantly between different types, models, and price points.
The word laptop, modeled after the term desktop (as in desktop computer), refers to the fact that the computer can be practically placed on the user's lap; while the word notebook refers to most laptops sharing a form factor with paper notebooks. , in American English, the terms laptop and notebook are used interchangeably; in other dialects of English, one or the other may be preferred. The term notebook originally referred to a type of portable computer that was smaller and lighter than mainstream laptops of the time, but has since come to mean the same thing and no longer refers to any specific size.
Design elements, form factors, and construction can also vary significantly between models depending on the intended use. Examples of specialized models of laptops include 2-in-1 laptops, with keyboards that either be detached or pivoted out of view from the display (often marketed having a "laptop mode"), and rugged laptops, for use in construction or military applications. Portable computers, which later developed into modern laptops, were originally considered to be a small niche market, mostly for specialized field applications, such as in the military, for accountants, or travelling sales representatives. As portable computers evolved into modern laptops, they became widely used for a variety of purposes.
History
The history of the laptop follows closely behind the development of the personal computer itself. A "personal, portable information manipulator" was imagined by Alan Kay at Xerox PARC in 1968, and described in his 1972 paper as the "Dynabook". The IBM Special Computer APL Machine Portable (SCAMP) was demonstrated in 1973. This prototype was based on the IBM PALM processor. The IBM 5100, the first commercially available portable computer, appeared in September 1975, and was based on the SCAMP prototype.
As 8-bit CPU machines became widely accepted, the number of portables increased rapidly. The first "laptop-sized notebook computer" was the Epson HX-20, invented (patented) by Suwa Seikosha's Yukio Yokozawa in July 1980, introduced at the COMDEX computer show in Las Vegas by Japanese company Seiko Epson in 1981, and released in July 1982. It had an LCD screen, a rechargeable battery, and a calculator-size printer, in a chassis, the size of an A4 notebook. It was described as a "laptop" and "notebook" computer in its patent.
Both Tandy/RadioShack and Hewlett-Packard (HP) also produced portable computers of varying designs during this period. The first laptops using the flip form factor appeared in the early 1980s. The Dulmont Magnum was released in Australia in 1981–82, but was not marketed internationally until 1984–85. The GRiD Compass 1101, released in 1982, was used at NASA and by the military, among others. The Sharp PC-5000, the Ampere WS-1, and Gavilan SC were released between 1983 and 1985. The Toshiba T1100 won acceptance by PC experts and the mass market as a way to have PC portability.
From 1983 onward, several new input techniques were developed and included in laptops, including the touch pad (Gavilan SC, 1983), the pointing stick (IBM ThinkPad 700, 1992), and handwriting recognition (Linus Write-Top, 1987). Some CPUs, such as the 1990 Intel i386SL, were designed to use minimum power to increase the battery life of portable computers and were supported by dynamic power management features such as Intel SpeedStep and AMD PowerNow! in some designs.
Some laptops in the 1980s using red plasma displays could only be used when connected to AC power, and had a built in power supply.
The development of memory cards was driven in the 1980s by the need for a floppy-disk-drive alternative, having lower power consumption, less weight, and reduced volume in laptops. The Personal Computer Memory Card International Association (PCMCIA) was an industry association created in 1989 to promote a standard for memory cards in PCs. The specification for PCMCIA type I cards, later renamed PC Cards, was first released in 1990.
Displays reached 640x480 (VGA) resolution by 1988 (Compaq SLT/286), and color screens started becoming a common upgrade in 1991, with increases in resolution and screen size occurring frequently until the introduction of 17" screen laptops in 2003. Hard drives started to be used in portables, encouraged by the introduction of 3.5" drives in the late 1980s, and became common in laptops starting with the introduction of 2.5" and smaller drives around 1990; capacities have typically lagged behind those of physically larger desktop drives.
Optical disc drives became common in full-size laptops around 1997: initially, CD-ROM drives, supplanted by CD-R, then DVD, then Blu-ray drives with writing capability. Starting around 2011, the trend shifted against internal optical drives, and as of 2022, they have largely disappeared, though are still readily available as external peripherals.
Resolutions of laptop webcams are 720p (HD), or 480p in lower-end laptops. The earliest-known laptops with 1080p (Full HD) webcams, like the Samsung 700G7C, were released in the early 2010s.
Etymology
The terms laptop and notebook trace their origins to the early 1980s, coined to describe portable computers in a size class smaller than the mainstream units (so-called "luggables") but larger than pocket computers. The etymologist William Safire traced the origin of laptop to some time before 1984; the earliest attestation of laptop found by the Oxford English Dictionary dates to 1983. The word is modeled after the term desktop, as in desktop computer. Notebook, meanwhile, emerged earlier in 1982 to describe Epson's HX-20 portable, whose dimensions roughly correspond to a letter-sized pad of paper. Notebooks emerged as their own separate market from laptops with the release of the NEC UltraLite in 1988.
Notebooks and laptops continued to occupy distinct market segments into the mid-1990s, but ergonomic considerations and customer preference for larger screens soon led to notebooks converging with laptops in the late 1990s. Now, the terms laptop and notebook are synonymous, with laptop being the more common term in most English-speaking territories.
Types of laptops
Since the 1970s introduction of portable computers, their forms have changed significantly, resulting in a variety of visually and technologically differing subclasses. Excepting distinct legal trademark around terms (notably Ultrabook), hard distinctions between these classes were rare, and their usage has varied over time and between sources. Since the late 2010s, more specific terms have become less commonly used, with sizes distinguished largely by the size of the screen.
Smaller and larger laptops
There were in the past a number of marketing categories for smaller and larger laptop computers; these included "notebook" and "subnotebook" models, low cost "netbooks", and "ultra-mobile PCs" where the size class overlapped with devices like smartphone and handheld tablets, and "Desktop replacement" laptops for machines notably larger and heavier than typical to operate more powerful processors or graphics hardware. All of these terms have fallen out of favor as the size of mainstream laptops has gone down and their capabilities have gone up; except for niche models, laptop sizes tend to be distinguished by the size of the screen, and for more powerful models, by any specialized purpose the machine is intended for, such as a "gaming laptop" or a "mobile workstation" for professional use.
Convertible, hybrid, 2-in-1
The latest trend of technological convergence in the portable computer industry spawned a broad range of devices, which combined features of several previously separate device types. The hybrids, convertibles, and 2-in-1s emerged as crossover devices, which share traits of both tablets and laptops. All such devices have a touchscreen display designed to allow users to work in a tablet mode, using either multi-touch gestures or a stylus/digital pen.
Convertibles are devices with the ability to conceal a hardware keyboard. Keyboards on such devices can be flipped, rotated, or slid behind the back of the chassis, thus transforming from a laptop into a tablet. Hybrids have a keyboard detachment mechanism, and due to this feature, all critical components are situated in the part with the display. 2-in-1s can have a hybrid or a convertible form, often dubbed 2-in-1 detachable and 2-in-1 convertibles respectively, but are distinguished by the ability to run a desktop OS, such as Windows 10. 2-in-1s are often marketed as laptop replacement tablets.
2-in-1s are often very thin, around , and light devices with a long battery life. 2-in-1s are distinguished from mainstream tablets as they feature an x86-architecture CPU (typically a low- or ultra-low-voltage model), such as the Intel Core i5, run a full-featured desktop OS like Windows 10, and have a number of typical laptop I/O ports, such as USB 3 and Mini DisplayPort.
2-in-1s are designed to be used not only as a media consumption device but also as valid desktop or laptop replacements, due to their ability to run desktop applications, such as Adobe Photoshop. It is possible to connect multiple peripheral devices, such as a mouse, keyboard, and several external displays to a modern 2-in-1.
Microsoft Surface Pro-series devices and Surface Book are examples of modern 2-in-1 detachable, whereas Lenovo Yoga-series computers are a variant of 2-in-1 convertibles. While the older Surface RT and Surface 2 have the same chassis design as the Surface Pro, their use of ARM processors and Windows RT do not classify them as 2-in-1s, but as hybrid tablets. Similarly, a number of hybrid laptops run a mobile operating system, such as Android. These include Asus's Transformer Pad devices, examples of hybrids with a detachable keyboard design, which do not fall in the category of 2-in-1s.
Rugged laptop
A rugged laptop is designed to reliably operate in harsh usage conditions such as strong vibrations, extreme temperatures, and wet or dusty environments. Rugged laptops are bulkier, heavier, and much more expensive than regular laptops, and thus are seldom seen in regular consumer use.
Hardware
The basic components of laptops function identically to their desktop counterparts. Traditionally they were miniaturized and adapted to mobile use, The design restrictions on power, size, and cooling of laptops limit the maximum performance of laptop parts compared to that of desktop components, although that difference has increasingly narrowed.
In general, laptop components are not intended to be replaceable or upgradable by the end-user, except for components that can be detached; in the past, batteries and optical drives were commonly exchangeable. Some laptops feature socketed processors with sockets such as the Socket G2, but many laptops use processors that are soldered to the motherboard. Many laptops come with RAM and storage that is soldered to the motherboard and cannot be easily replaced. This restriction is one of the major differences between laptops and desktop computers, because the large "tower" cases used in desktop computers are designed so that new motherboards, hard disks, sound cards, RAM, and other components can be added. Memory and storage can often be upgraded with some disassembly, but with the most compact laptops, there may be no upgradeable components at all.
The following sections summarize the differences and distinguishing features of laptop components in comparison to desktop personal computer parts.
Display
The typical laptop has a screen that, when unfolded, is upright to the user.
Screen technology
Laptop screens most commonly use liquid-crystal display (LCD) technology, although OLED panels have been used in some models. The display interfaces with the motherboard using the embedded DisplayPort protocol via the Low-voltage differential signaling (LVDS) 30 or 40 pin connector. Earlier laptops use the FPD-Link standard. The panels are mainly manufactured by AU Optronics, BOE Technology, LG Display or Samsung Display.
Surface finish
Externally, it can be a glossy or a matte (anti-glare) screen.
Sizes
In the past, there was a broader range of marketing terms (both formal and informal) to distinguish between different sizes of laptops. These included netbooks, subnotebooks, ultra-mobile PC, and desktop replacement computers; these are sometimes still used informally, although they are generally not used anymore in manufacturer marketing.
mainstream consumer laptops tend to come with 11", 13" or 15"-16" screens; 14" models are more popular among business machines. Larger and smaller models are available, but less common – there is no clear dividing line in minimum or maximum size. Machines small enough to be handheld (screens in the 6–8" range) can be marketed either as very small laptops or "handheld PCs", while the distinction between the largest laptops and "All-in-One" desktops is whether they fold for travel.
Resolution
Having a higher resolution display allows more items to fit onscreen at a time, improving the user's ability to multitask, although, at the higher resolutions on smaller screens, the resolution may only serve to display sharper graphics and text rather than increasing the usable area. Since the introduction of the MacBook Pro with Retina display in 2012, there has been an increase in the availability of "HiDPI" (or high pixel density) displays; this is generally considered to be anything higher than 1920 pixels wide. This has increasingly converged around 4K (3840-pixel-wide) resolutions.
External displays can be connected to most laptops, with most models supporting at least one. The use of technology such as USB4 (section Alternate Mode partner specifications). DisplayPort Alt Mode has been utilized to charge a laptop and provide display output over one USB-C Cable.
Refresh rates
Most laptop displays have a maximum refresh rate of 60 Hz. The Dell M17x and Samsung 700G7A, both released in 2011, were among the first laptops to feature a 120 Hz refresh rate, and more such laptops have appeared in the years since.
Central processing unit (CPU)
Laptop CPUs have advanced power-saving features and produce less heat than those intended for desktop use. Mainstream laptop CPUs made after 2018 have at least two processor cores, often four cores, and sometimes more, with 6 and 8 cores becoming more common.
For the low price and mainstream performance, there is no longer a significant performance difference between laptop and desktop CPUs, but at the high end, the fastest desktop CPUs still substantially outperform the fastest laptop processors, at the expense of massively higher power consumption and heat generation; the fastest laptop processors top out at 56 watts of heat, while the fastest desktop processors top out at 150 watts (and often need water cooling).
There has been a wide range of CPUs designed for laptops available from both Intel, AMD, and other manufacturers. On non-x86 architectures, Motorola and IBM produced the chips for the former PowerPC-based Apple laptops (iBook and PowerBook). Between around 2000 to 2014, most full-size laptops had socketed, replaceable CPUs; on thinner models, the CPU was soldered on the motherboard and was not replaceable or upgradable without replacing the motherboard. Since 2015, Intel has not offered new laptop CPU models with pins to be interchangeable, preferring ball grid array chip packages which have to be soldered; and as of 2021, only a few rare models using desktop parts.
In the past, some laptops have used a desktop processor instead of the laptop version and have had high-performance gains at the cost of greater weight, heat, and limited battery life; this is not unknown as of 2022, but since around 2010, the practice has been restricted to small-volume gaming models. Laptop CPUs are rarely able to be overclocked; most use locked processors. Even on gaming models where unlocked processors are available, the cooling system in most laptops is often very close to its limits and there is rarely headroom for an overclocking–related operating temperature increase.
Graphics processing unit (GPU)
On most laptops, the GPU is integrated into the CPU to conserve power and space. This was introduced by Intel with the Core i-series of mobile processors in 2010, followed by similar AMD APU processors in January 2011.
Before that, lower-end machines tended to use graphics processors integrated into the system chipset, while higher-end machines had a separate graphics processor. In the past, laptops lacking a separate graphics processor were limited in their utility for gaming and professional applications involving 3D graphics, but the capabilities of CPU-integrated graphics have converged with the low-end of dedicated graphics processors since the mid-2010s. For laptops possessing limited onboard graphics capability but sufficient I/O throughput, an external GPU (eGPU) can provide additional graphics power at the cost of physical space and portability.
Higher-end laptops intended for gaming or professional 3D work still come with dedicated (and in some cases even dual) graphics processors on the motherboard or as an internal expansion card. Since 2011, these almost always involve switchable graphics so that when there is no demand for the higher performance dedicated graphics processor, the more power-efficient integrated graphics processor will be used. Nvidia Optimus and AMD Hybrid Graphics are examples of this sort of system of switchable graphics.
Traditionally, the system RAM on laptops (as well as on desktop computers) was physically separate from the graphics memory used by the GPU. Apple's M series SoCs feature a unified pool of memory for both the system and the GPU; this approach can produce substantial efficiency gains for some applications but comes at the cost of eGPU support.
Memory
Since around the year 2000, most laptops have used SO-DIMM slots in which RAM is mounted, although, as of 2021, an increasing number of models use memory soldered to the motherboard, either alongside SO-DIMM slots or without any slots and soldering all memory to the motherboard. A new form factor, the CAMM module, is slated to fix the size and timing limitation. Before 2000, most laptops used proprietary memory modules if their memory was upgradable.
In the early 2010s, high end laptops such as the 2011 Samsung 700G7A have passed the 10 GB RAM barrier, featuring 16 GB of RAM.
When upgradeable, memory slots are sometimes accessible from the bottom of the laptop for ease of upgrading; in other cases, accessing them requires significant disassembly. Most laptops have two memory slots, although some will have only one, either for cost savings or because some amount of memory is soldered. Some high-end models have four slots; these are usually mobile engineering workstations, although a few high-end models intended for gaming do as well.
8 GB RAM is most common, with lower-end models occasionally having 4 GB. Higher-end laptops may come with 16 GB of RAM or more.
Internal storage
The earliest laptops most often used floppy disks for storage, although a few used either RAM disk or tape. By the late 1980s hard disk drives had become the standard form of storage.
Between 1990 and 2009, almost all laptops typically had a hard disk drive (HDD) for storage; since then, solid-state drives (SSD) have gradually come to replace hard drives in all but some inexpensive consumer models. Solid-state drives are faster and more power-efficient, as well as eliminating the hazard of drive and data corruption caused by a laptop's physical impacts, as they use no mechanical parts such as a rotational platter. In many cases, they are more compact as well. Initially, in the late 2000s, SSDs were substantially more expensive than HDDs, but prices on smaller capacity (under 1 terabyte) drives have converged; larger capacity drives remain more expensive than comparable-sized HDDs.
Since around 1990, where a hard drive is present it will typically be a 2.5-inch drive; some very compact laptops support even smaller 1.8-inch HDDs, and a very small number used 1" Microdrives. Some SSDs are built to match the size/shape of a laptop hard drive, but increasingly they have been replaced with smaller mSATA or M.2 cards. SSDs using the newer and much faster NVM Express standard for connecting are only available as cards.
many laptops no longer contain space for a 2.5" drive, accepting only M.2 cards; a few of the smallest have storage soldered to the motherboard. For those that can, they can typically contain a single 2.5-inch drive, but a small number of laptops with a screen wider than 15 inches can house two drives.
A variety of external HDDs or NAS data storage servers with support of RAID technology can be attached to virtually any laptop over such interfaces as USB, FireWire, eSATA, or Thunderbolt, or over a wired or wireless network to further increase space for the storage of data. Many laptops also incorporate a SD or microSD card slot. This enables users to download digital pictures from an SD card onto a laptop, thus enabling them to delete the SD card's contents to free up space for taking new pictures.
Removable media drive
Optical disc drives capable of playing CD-ROMs, compact discs (CD), DVDs, and in some cases, Blu-ray discs (BD), were nearly universal on full-sized models between the mid-1990s and the early 2010s. drives are uncommon in compact or premium laptops; they remain available in some bulkier models, but the trend towards thinner and lighter machines is gradually eliminating these drives and players – when needed they can be connected via USB instead.
Speaker
Laptops usually have built-in speakers and built-in microphones. However, integrated speakers may be small and of restricted sound quality to conserve space.
Inputs
An alphanumeric keyboard is used to enter text, data, and other commands (e.g., function keys). A touchpad (also called a trackpad), a pointing stick, or both, are used to control the position of the cursor on the screen, and an integrated keyboard is used for typing. Some touchpads have buttons separate from the touch surface, while others share the surface. A quick double-tap is typically registered as a click, and operating systems may recognize multi-finger touch gestures.
An external keyboard and mouse may be connected using a USB port or wirelessly, via Bluetooth or similar technology. Some laptops have multitouch touchscreen displays, either available as an option or standard. Most laptops have webcams and microphones, which can be used to communicate with other people with both moving images and sound, via web conferencing or video-calling software.
Laptops typically have USB ports and a combined headphone/microphone jack, for use with headphones, a combined headset, or an external mic. Many laptops have a card reader for reading digital camera SD cards.
Input/output (I/O) ports
On a typical laptop, there are several USB ports; if they use only the older USB connectors instead of USB-C, they will typically have an external monitor port (VGA, DVI, HDMI or Mini DisplayPort or occasionally more than one), an audio in/out port (often in form of a single socket) is common. It is possible to connect up to three external displays to a 2014-era laptop via a single Mini DisplayPort, using multi-stream transport technology.
Apple, in a 2015 version of its MacBook, transitioned from a number of different I/O ports to a single USB-C port. This port can be used both for charging and connecting a variety of devices through the use of aftermarket adapters. Apple has since transitioned back to using a number of different ports. Google, with its updated version of Chromebook Pixel, shows a similar transition trend towards USB-C, although keeping older USB Type-A ports for a better compatibility with older devices. Although being common until the end of the 2000s decade, Ethernet network port are rarely found on modern laptops, due to widespread use of wireless networking, such as Wi-Fi. Legacy ports such as a PS/2 keyboard/mouse port, serial port, parallel port, or FireWire are provided on some models, but they are increasingly rare. On Apple's systems, and on a handful of other laptops, there are also Thunderbolt ports, but Thunderbolt 3 uses USB-C. Laptops typically have a headphone jack, so that the user can connect headphones or amplified speaker systems for listening to music or other audio.
Expansion cards
In the past, a PC Card (formerly PCMCIA) or ExpressCard slot for expansion was often present on laptops to allow adding and removing functionality, even when the laptop is powered on; these are becoming increasingly rare since the introduction of USB 3.0. Some internal subsystems such as Ethernet, Wi-Fi, or a wireless cellular modem can be implemented as replaceable internal expansion cards, usually accessible under an access cover on the bottom of the laptop. The standard for such cards is PCI Express, which comes in both mini and even smaller M.2 sizes. In newer laptops, it is not uncommon to also see Micro SATA (mSATA) functionality on PCI Express Mini or M.2 card slots allowing the use of those slots for SATA-based solid-state drives.
Mobile PCI Express Module (MXM) is a type of expansion card that is used for graphics cards.
Battery and power supply
Since the late 1990s, laptops have typically used lithium ion or lithium polymer batteries, These replaced the older nickel metal-hydride typically used in the 1990s, and nickel–cadmium batteries used in most of the earliest laptops. A few of the oldest laptops used non-rechargeable batteries, or lead–acid batteries.
Battery life is highly variable by model and workload and can range from one hour to nearly a day. A battery's performance gradually decreases over time; a noticeable reduction in capacity is typically evident after two to three years of regular use, depending on the charging and discharging pattern and the design of the battery. Innovations in laptops and batteries have seen situations in which the battery can provide up to 24 hours of continued operation, assuming average power consumption levels. An example is the HP EliteBook 6930p when used with its ultra-capacity battery.
Laptops with removable batteries may support larger replacement batteries with extended capacity.
A laptop's battery is charged using an external power supply, which is plugged into a wall outlet. The power supply outputs a DC voltage typically in the range of 7.2—24 volts. The power supply is usually external and connected to the laptop through a DC connector cable. In most cases, it can charge the battery and power the laptop simultaneously. When the battery is fully charged, the laptop continues to run on power supplied by the external power supply, avoiding battery use. If the used power supply is not strong enough to power computing components and charge the battery simultaneously, the battery may charge in a shorter period of time if the laptop is turned off or sleeping. The charger typically adds about to the overall transporting weight of a laptop, although some models are substantially heavier or lighter. Most 2016-era laptops use a smart battery, a rechargeable battery pack with a built-in battery management system (BMS). The smart battery can internally measure voltage and current, and deduce charge level and State of Health (SoH) parameters, indicating the state of the cells.
Power connectors
Historically, DC connectors, typically cylindrical/barrel-shaped coaxial power connectors have been used in laptops. Some vendors such as Lenovo made intermittent use of a rectangular connector.
Some connector heads feature a center pin to allow the end device to determine the power supply type by measuring the resistance between it and the connector's negative pole (outer surface). Vendors may block charging if a power supply is not recognized as the original part, which could deny the legitimate use of universal third-party chargers.
With the advent of USB-C, portable electronics made increasing use of it for both power delivery and data transfer. Its support for 20 V (common laptop power supply voltage) and 5 A typically suffices for low to mid-end laptops, but some with higher power demands such as gaming laptops depend on dedicated DC connectors to handle currents beyond 5 A without risking overheating, some even above 10 A. Additionally, dedicated DC connectors are more durable and less prone to wear and tear from frequent reconnection, as their design is less delicate.
Cooling
Waste heat from the operation is difficult to remove in the compact internal space of a laptop. The earliest laptops used passive cooling; this gave way to heat sinks placed directly on the components to be cooled, but when these hot components are deep inside the device, a large space-wasting air duct is needed to exhaust the heat. Modern laptops instead rely on heat pipes to rapidly move waste heat towards the edges of the device, to allow for a much smaller and compact fan and heat sink cooling system. Waste heat is usually exhausted away from the device operator towards the rear or sides of the device. Multiple air intake paths are used since some intakes can be blocked, such as when the device is placed on a soft conforming surface like a chair cushion. Secondary device temperature monitoring may reduce performance or trigger an emergency shutdown if it is unable to dissipate heat, such as if the laptop were to be left running and placed inside a carrying case. Aftermarket cooling pads with external fans can be used with laptops to reduce operating temperatures.
Docking station
A docking station (sometimes referred to simply as a dock) is a laptop accessory that contains multiple ports and in some cases expansion slots or bays for fixed or removable drives. A laptop connects and disconnects to a docking station, typically through a single large proprietary connector. A docking station is an especially popular laptop accessory in a corporate computing environment, due to the possibility of a docking station transforming a laptop into a full-featured desktop replacement, yet allowing for its easy release. This ability can be advantageous to "road warrior" employees who have to travel frequently for work, and yet who also come into the office. If more ports are needed, or their position on a laptop is inconvenient, one can use a cheaper passive device known as a port replicator. These devices mate to the connectors on the laptop, such as through USB or FireWire.
Charging trolleys
Laptop charging trolleys, also known as laptop trolleys or laptop carts, are mobile storage containers to charge multiple laptops, netbooks, and tablet computers at the same time. The trolleys are used in schools that have replaced their traditional static computer labs suites of desktop equipped with "tower" computers, but do not have enough plug sockets in an individual classroom to charge all of the devices. The trolleys can be wheeled between rooms and classrooms so that all students and teachers in a particular building can access fully charged IT equipment.
Laptop charging trolleys are also used to deter and protect against opportunistic and organized theft. Schools, especially those with open plan designs, are often prime targets for thieves who steal high-value items. Laptops, netbooks, and tablets are among the highest–value portable items in a school. Moreover, laptops can easily be concealed under clothing and stolen from buildings. Many types of laptop–charging trolleys are designed and constructed to protect against theft. They are generally made out of steel, and the laptops remain locked up while not in use. Although the trolleys can be moved between areas from one classroom to another, they can often be mounted or locked to the floor, support pillars, or walls to prevent thieves from stealing the laptops, especially overnight.
Solar panels
In some laptops, solar panels are able to generate enough solar power for the laptop to operate. The One Laptop Per Child Initiative released the OLPC XO-1 laptop which was tested and successfully operated by use of solar panels. They were designing an OLPC XO-3 laptop with these features. The OLPC XO-3 was planned to operate with 2 watts of electricity. Samsung has also designed the NC215S solar–powered notebook that was planned to be sold commercially in the U.S. market.
Accessories
A common accessory for laptops is a laptop sleeve, laptop skin, or laptop case, which provides a degree of protection from scratches. Sleeves, which are distinguished by being relatively thin and flexible, are most commonly made of neoprene, with sturdier ones made of low-resilience polyurethane. Some laptop sleeves are wrapped in ballistic nylon to provide some measure of waterproofing. Bulkier and sturdier cases can be made of metal with polyurethane padding inside and may have locks for added security. Metal, padded cases also offer protection against impacts and drops. Another common accessory is a laptop cooler, a device that helps lower the internal temperature of the laptop either actively or passively. A common active method involves using electric fans to draw heat away from the laptop, while a passive method might involve propping the laptop up on some type of pad so it can receive more airflow. Some stores sell laptop pads that enable a reclining person on a bed to use a laptop.
Modularity
Some of the components of earlier models of laptops can easily be replaced without opening completely its bottom part, such as the keyboard, battery, hard disk, memory modules, and CPU cooling fan.
Some of the components of recent models of laptops reside inside. Replacing most of its components, such as the keyboard, battery, hard disk, memory modules, CPU cooling fan, etc., requires the removal of either the top or bottom part, the removal of the motherboard, and returning them.
In some types, solder and glue are used to mount components such as RAM, storage, and batteries, making repairs additionally difficult.
Obsolete features
Features that certain early models of laptops used to have that are not available in more recent models include:
Reset ("cold restart") button in a hole (needed a thin metal tool to press)
Instant power off button in a hole (needed a thin metal tool to press)
Integrated charger or power adapter inside the laptop
Dedicated Media buttons (Internet, Volume, Play, Pause, Next, Previous)
Floppy disk drive
Serial port
Parallel port
Modem
IEEE 1394 port
Docking port
Shared PS/2 input device port
IrDA
S-video port
S/PDIF audio port
PC Card / PCMCIA slot
ExpressCard slot
CD/DVD Drives (starting with 2013 models)
VGA port (starting with 2013 models)
Characteristics
Advantages over desktop computers
Portability - Laptops are highly portable compared to desktop PCs. Physical portability allows a laptop to be used in many places—not only at home and the office but also during commuting and flights, in coffee shops, in lecture halls and libraries, at clients' locations or a meeting room, etc. Within a home, portability enables laptop users to move their devices from room to room. Portability offers several distinct advantages:
Productivity: Using a laptop in places where a desktop PC cannot be used can help employees and students to increase their productivity on work or school tasks, such as an office worker reading their work e-mails during an hour-long commute by train, or a student doing their homework at the university coffee shop during a break between lectures, for example.
Up-to-date information: Using a single laptop prevents fragmentation of files across multiple PCs as the files exist in a single location and are always up-to-date.
Connectivity: A key advantage of laptops is that they almost always have integrated connectivity features such as Wi-Fi and Bluetooth, and sometimes connection to cellular networks either through native integration or use of a hotspot. Wi-Fi networks and laptop programs are especially widespread at university campuses.
Other advantages of laptops:
Size: Laptops are smaller than desktop PCs. This is beneficial when space is at a premium, for example in small apartments and student dorms. When not in use, a laptop can be closed and put away in a desk drawer.
Low power consumption: Laptops are several times more power-efficient than desktops. A typical laptop uses 10–100 W, compared to 200–800W for desktops. This could be particularly beneficial for large businesses, which run hundreds of personal computers thus economies of scale, and homes where there is a computer running 24/7 (such as a home media server, print server, etc.).
Quiet: Laptops are typically much quieter than desktops, due both to the components (often silent solid-state drives replacing hard drives) and to less heat production leading to the use of fewer, sometimes no cooling fans. The latter has given rise to laptops that have no moving parts, resulting in complete silence during use.
Battery: a charged laptop can continue to be used in case of a power outage and is not affected by short power interruptions and blackouts, an issue that is present with desktop PCs.
All-in-One: designed to be portable, most modern laptops have all components integrated into the chassis. For desktops (excluding all-in-ones) this is usually divided into the desktop "tower" (the unit with the CPU, hard drive, power supply, etc.), keyboard, mouse, display screen, and optional peripherals such as speakers.
Disadvantages
Compared to desktop PCs, laptops have disadvantages in the following areas:
Performance The performance of laptops is often worse than comparably priced desktops. The upper limits of performance of laptops remain lower than desktops, due to mostly practical reasons, such as decreased battery life, increased size and heat, etc.
Upgradeability The upgradeability of laptops is limited compared to tower desktops, due to technical and economic reasons. In general, hard drives and memory can be upgraded easily. Due to the integrated nature of laptops, however, the motherboard, CPU, and graphics, are seldom officially upgradeable. Some efforts towards industry standard parts and layouts have been attempted, such as Common Building Block, but the industry remains largely proprietary and fragmented. There is no industry-wide standard form factor for laptops; Moreover, starting with 2013 models, laptops have become increasingly integrated (soldered) with the motherboard for most of its components (CPU, SSD, RAM, etc.) to reduce size and upgradeability prospects.
Durability Laptops are less durable than desktops/PCs. However, the durability of the laptop depends on the user if proper maintenance is done then the laptop can work longer.Because of their portability, laptops are subject to more wear and physical damage than desktops, additionally hindered by their integrated nature. A liquid spill onto the keyboard, while a minor issue with a desktop system, can damage the internals of a laptop and destroy the computer, resulting in a costly repair or entire replacement of laptops. One study found that a laptop is three times more likely to break during the first year of use than a desktop. To maintain a laptop, it is recommended to clean it every three months for dirt, debris, dust, and food particles. Most cleaning kits consist of a lint-free or microfiber cloth for the screen and keyboard, compressed air for getting dust out of the cooling fan, and a cleaning solution. Harsh chemicals such as bleach should not be used to clean a laptop, as they can damage it.
Heating and cooling Laptops rely on extremely compact cooling systems involving a fan and heat sink that can fail from blockage caused by accumulated airborne dust and debris. Most laptops do not have any type of removable dust collection filter over the air intake for these cooling systems, resulting in a system that gradually conducts more heat and noise as the years pass. In some cases, the laptop starts to overheat even at idle load levels. This dust is usually stuck inside where the fan and heat sink meet, where it can not be removed by a casual cleaning and vacuuming. Most of the time, compressed air can dislodge the dust and debris but may not entirely remove it. After the device is turned on, the loose debris is reaccumulated into the cooling system by the fans. Complete disassembly is usually required to clean the laptop entirely. However, preventative maintenance such as regular cleaning of the heat sink via compressed air can prevent dust build-up on the heat sink. Many laptops are difficult to disassemble by the average user and contain components that are sensitive to electrostatic discharge (ESD).
Battery life Battery life is limited because the capacity drops with time, eventually warranting replacement after as little as 2–3 years. A new battery typically stores enough energy to run the laptop for five to six hours or more, depending on usage and the battery size. The battery is often easily replaceable and a higher capacity model may be obtained for longer charging and discharging time. Some laptops do not have the usual removable battery and have to be brought to the service center of their manufacturer or a third-party laptop service center to have their battery replaced. Replacement batteries can also be expensive, depending on the availability of the parts. Desktop PCs do not face similar problems since they are reliant on long lasting power supplies.
Security and privacyBecause they are valuable, commonly used, portable, and easy to hide in a backpack or other type of bag, laptops are often stolen. Every day, over 1,600 laptops go missing from U.S. airports. The cost of stolen business or personal data, and of the resulting problems (identity theft, credit card fraud, breach of privacy), can be many times the value of the stolen laptop itself. Consequently, the physical protection of laptops and the safeguarding of data contained in them are both of great importance. Some laptops, primarily professional and educational devices, have a Kensington security slot, which can be used to tether them with a security cable and lock. In addition, modern operating systems have features such as Activation Lock or similar that prevents the use of the device without credentials. some laptops also have additional security elements added, including biometric security components such as Windows Hello or Touch ID.Software such as GadgetTrak and Find My Mac have been engineered to help people locate and recover their stolen laptops in the event of theft. Setting one's laptop with a password on its firmware (protection against going to firmware setup or booting), internal HDD/SSD (protection against accessing it and loading an operating system on it afterward), and every user account of the operating system are additional security measures that a user should do. Fewer than 5% of lost or stolen laptops are recovered by the companies that own them, however, that number may decrease due to a variety of companies and software solutions specializing in laptop recovery. In the 2010s, the common availability of webcams on laptops raised privacy concerns. In Robbins v. Lower Merion School District (Eastern District of Pennsylvania 2010), school-issued laptops loaded with special software enabled staff from two high schools to take secret webcam shots of students at home, via their students' laptops.
Ergonomics and health effects
Wrists Prolonged use of laptops can cause repetitive strain injury because of their small, flat keyboard and trackpad pointing devices. Usage of separate, external ergonomic keyboards and pointing devices is recommended to prevent injury when working for long periods of time; they can be connected to a laptop easily by USB, Bluetooth or via a docking station. Some health standards require ergonomic keyboards at workplaces.
Neck and spine A laptop's integrated screen often requires users to lean over for a better view, which can cause neck or spinal injuries. A larger and higher-quality external screen can be connected to almost any laptop to alleviate this and to provide additional screen space for more productive work. Another solution is to use a computer stand.
Possible effect on fertility A study by State University of New York researchers found that heat generated from laptops can increase the temperature of the lap of male users when balancing the computer on their lap, potentially putting sperm count at risk. The study, which included roughly two dozen men between the ages of 21 and 35, found that the sitting position required to balance a laptop can increase scrotum temperature by as much as . However, further research is needed to determine whether this directly affects male sterility. A later 2010 study of 29 males published in Fertility and Sterility found that men who kept their laptops on their laps experienced scrotal hyperthermia (overheating) in which their scrotal temperatures increased by up to . The resulting heat increase, which could not be offset by a laptop cushion, may increase male infertility. A common practical solution to this problem is to place the laptop on a table or desk or to use a book or pillow between the body and the laptop. Another solution is to obtain a cooling unit for the laptop. These are usually USB powered and consist of a hard thin plastic case housing one, two, or three cooling fans – with the entire assembly designed to sit under the laptop in question – which results in the laptop remaining cool to the touch, and greatly reduces laptop heat buildup.
ThighsHeat generated from using a laptop on the lap can also cause skin discoloration on the thighs known as "toasted skin syndrome".
Sales
Manufacturers
There are many laptop brands and manufacturers. Several major brands that offer notebooks in various classes are listed in the adjacent box. The major brands usually offer good service and support, including well-executed documentation and driver downloads that remain available for many years after a particular laptop model is no longer produced. Capitalizing on service, support, and brand image, laptops from major brands are more expensive than laptops from smaller brands and ODMs. Some brands specialize in a particular class of laptops, such as gaming laptops (Alienware), high-performance laptops (HP Envy), netbooks (EeePC) and laptops for children (OLPC).
Many brands, including the major ones, do not design and do not manufacture their laptops. Instead, a small number of Original Design Manufacturers (ODMs) design new models of laptops, and the brands choose the models to be included in their lineup. In 2006, 7 major ODMs manufactured 7 of every 10 laptops in the world, with the largest one (Quanta Computer) having 30% of the world market share. Therefore, identical models are available both from a major label and from a low-profile ODM in-house brand.
Historic market share
As of 1992–1993, Toshiba ranked as the global leading vendor in the notebook computer market. In the United States meanwhile, Apple led the market followed by Compaq. In the year 1993, global revenue for the laptop market was led by Compaq, followed by Toshiba, Apple, NEC and IBM, altogether accounting for over 53% of global revenue.
In the United States, the top three vendors for notebooks in market share as of 1996 were: Toshiba, followed by Compaq, and followed by IBM.
As of 1999, Toshiba ranked first in worldwide laptop sales followed by IBM, Compaq, and Dell. Toshiba led the market with a share of 18.6%.
In the first quarter of 2002 in the United States market, Dell controlled 25.2% in the notebook space, well ahead of Toshiba (13.6%) and Compaq (11.7%), the latter of which had been acquired by Hewlett-Packard (HP). At fourth and fifth place were Sony and IBM.
In Europe, the Middle East and Africa (EMEA) territories, Acer was the largest vendor of laptops, in 2004–2005, having overtaken HP and IBM there.
In the year 2005 according to IDC, Dell was the top global vendor of notebooks with a market share of 17.29%, followed by: HP (15.7%), Toshiba (10.96%), Acer (10.15%) and Lenovo (8.23%); Lenovo had acquired IBM that same year. The remaining of the top ten was made up of Fujitsu Siemens, Sony, NEC, Apple and Asus.
In the first quarter of 2010, the largest vendor of portable computers, including netbooks, was either HP or Acer, depending on data source. Both had shipped approximately 9 million units each. Dell, Toshiba, Asus and Lenovo followed, each with approximate sales of 5 to 6 million each. Apple, Samsung and Sony sold under 2 million each.
As of the third quarter of 2020, HP was cited as the leading vendor for notebook computers closely followed by Lenovo, both with a share of 23.6% each. They were followed by Dell (13.7%), Apple (9.7%) and Acer (7.9%).
Adoption by users
Battery-powered portable computers had just 2% worldwide market share in 1986. However, laptops have become increasingly popular, both for business and personal use. The third quarter of 2008 was the first time when worldwide notebook PC shipments exceeded desktops, with 38.6 million units versus 38.5 million units. In 2023, it was estimated that 166 million laptops were sold, and in the first quarter of 2024, around 64% of personal computers sold were laptops or detachable tablets. Due to the advent of tablets and affordable laptops, many computer users now have laptops due to the convenience offered by the device.
Price
Before 2008, laptops were very expensive. In May 2005, the average notebook sold for while desktops sold for an average of . Around 2008, however, prices of laptops decreased substantially due to low-cost netbooks, drawing an average at U.S. retail stores in August 2008. Starting with the 2010s, laptops have decreased substantially in price at the low end due to inexpensive and low power Arm processors, less demanding operating systems such as ChromeOS, and SoC's. , a new laptop can be obtained for .
Disposal
The list of materials that go into a laptop computer is long, and many of the substances used, such as beryllium, lead, chromium, and mercury compounds, are toxic or carcinogenic to humans. Although these toxins are relatively harmless when the laptop is in use, concerns that discarded laptops cause a serious health and environmental risks when improperly discarded have arisen. The Waste Electrical and Electronic Equipment Directive (WEEE Directive) in Europe specified that all laptop computers must be recycled by law. Similarly, the U.S. Environmental Protection Agency (EPA) has outlawed landfill dumping or the incinerating of discarded laptop computers.
Most laptop computers begin the recycling process with a method known as Demanufacturing, which involves the physical separation of the components of the laptop. These components are then either grouped into materials (e.g. plastic, metal and glass) for recycling or more complex items that require more advanced materials separation (e.g.) circuit boards, hard drives and batteries.
Corporate laptop recycling can require an additional process known as data destruction. The data destruction process ensures that all information or data that has been stored on a laptop hard drive can never be retrieved again. Below is an overview of some of the data protection and environmental laws and regulations applicable for laptop recycling data destruction:
Data Protection Act 1998 (DPA)
EU Privacy Directive (Due 2016)
Financial Conduct Authority
Sarbanes-Oxley Act
PCI-DSS Data Security Standard
Waste, Electronic & Electrical Equipment Directive (WEEE)
Basel Convention
Bank Secrecy Act (BSA)
FACTA Sarbanes-Oxley Act
FDA Security Regulations (21 C.F.R. part 11)
Gramm-Leach-Bliley Act (GLBA)
HIPAA (Health Insurance Portability and Accountability Act)
NIST SP 800–53
Add NIST SP 800–171
Identity Theft and Assumption Deterrence Act
Patriot Act of 2002
PCI Data Security Standard
US Safe Harbor Provisions
Various state laws
6/3 JAN
Gramm-leach-Bliley Act
DCID
Extreme use
The ruggedized Grid Compass computer was used since the early days of the Space Shuttle program. The first commercial laptop used in space was a Macintosh portable in 1990 on Space Shuttle mission STS-41 and again in 1991 aboard STS-43. Apple and other laptop computers continue to be flown aboard crewed spaceflights, though the only long-duration flight certified computer for the International Space Station is the ThinkPad. As of 2011, over 100 ThinkPads were aboard the ISS. Laptops used aboard the International Space Station and other spaceflights are generally the same ones that can be purchased by the general public but needed modifications are made to allow them to be used safely and effectively in a weightless environment such as updating the cooling systems to function without relying on hot air rising and accommodation for the lower cabin air pressure. Laptops operating in harsh usage environments and conditions, such as strong vibrations, extreme temperatures, and wet or dusty conditions differ from those used in space in that they are custom designed for the task and do not use commercial off-the-shelf hardware.
| Technology | Computer hardware | null |
198608 | https://en.wikipedia.org/wiki/Molecular%20dynamics | Molecular dynamics | Molecular dynamics (MD) is a computer simulation method for analyzing the physical movements of atoms and molecules. The atoms and molecules are allowed to interact for a fixed period of time, giving a view of the dynamic "evolution" of the system. In the most common version, the trajectories of atoms and molecules are determined by numerically solving Newton's equations of motion for a system of interacting particles, where forces between the particles and their potential energies are often calculated using interatomic potentials or molecular mechanical force fields. The method is applied mostly in chemical physics, materials science, and biophysics.
Because molecular systems typically consist of a vast number of particles, it is impossible to determine the properties of such complex systems analytically; MD simulation circumvents this problem by using numerical methods. However, long MD simulations are mathematically ill-conditioned, generating cumulative errors in numerical integration that can be minimized with proper selection of algorithms and parameters, but not eliminated.
For systems that obey the ergodic hypothesis, the evolution of one molecular dynamics simulation may be used to determine the macroscopic thermodynamic properties of the system: the time averages of an ergodic system correspond to microcanonical ensemble averages. MD has also been termed "statistical mechanics by numbers" and "Laplace's vision of Newtonian mechanics" of predicting the future by animating nature's forces and allowing insight into molecular motion on an atomic scale.
History
MD was originally developed in the early 1950s, following earlier successes with Monte Carlo simulationswhich themselves date back to the eighteenth century, in the Buffon's needle problem for examplebut was popularized for statistical mechanics at Los Alamos National Laboratory by Marshall Rosenbluth and Nicholas Metropolis in what is known today as the Metropolis–Hastings algorithm. Interest in the time evolution of N-body systems dates much earlier to the seventeenth century, beginning with Isaac Newton, and continued into the following century largely with a focus on celestial mechanics and issues such as the stability of the solar system. Many of the numerical methods used today were developed during this time period, which predates the use of computers; for example, the most common integration algorithm used today, the Verlet integration algorithm, was used as early as 1791 by Jean Baptiste Joseph Delambre. Numerical calculations with these algorithms can be considered to be MD done "by hand".
As early as 1941, integration of the many-body equations of motion was carried out with analog computers. Some undertook the labor-intensive work of modeling atomic motion by constructing physical models, e.g., using macroscopic spheres. The aim was to arrange them in such a way as to replicate the structure of a liquid and use this to examine its behavior. J.D. Bernal describes this process in 1962, writing:... I took a number of rubber balls and stuck them together with rods of a selection of different lengths ranging from 2.75 to 4 inches. I tried to do this in the first place as casually as possible, working in my own office, being interrupted every five minutes or so and not remembering what I had done before the interruption.Following the discovery of microscopic particles and the development of computers, interest expanded beyond the proving ground of gravitational systems to the statistical properties of matter. In an attempt to understand the origin of irreversibility, Enrico Fermi proposed in 1953, and published in 1955, the use of the early computer MANIAC I, also at Los Alamos National Laboratory, to solve the time evolution of the equations of motion for a many-body system subject to several choices of force laws. Today, this seminal work is known as the Fermi–Pasta–Ulam–Tsingou problem. The time evolution of the energy from the original work is shown in the figure to the right.
In 1957, Berni Alder and Thomas Wainwright used an IBM 704 computer to simulate perfectly elastic collisions between hard spheres. In 1960, in perhaps the first realistic simulation of matter, J.B. Gibson et al. simulated radiation damage of solid copper by using a Born–Mayer type of repulsive interaction along with a cohesive surface force. In 1964, Aneesur Rahman published simulations of liquid argon that used a Lennard-Jones potential; calculations of system properties, such as the coefficient of self-diffusion, compared well with experimental data. Today, the Lennard-Jones potential is still one of the most frequently used intermolecular potentials. It is used for describing simple substances (a.k.a. Lennard-Jonesium) for conceptual and model studies and as a building block in many force fields of real substances.
Areas of application and limits
First used in theoretical physics, the molecular dynamics method gained popularity in materials science soon afterward, and since the 1970s it has also been commonly used in biochemistry and biophysics. MD is frequently used to refine 3-dimensional structures of proteins and other macromolecules based on experimental constraints from X-ray crystallography or NMR spectroscopy. In physics, MD is used to examine the dynamics of atomic-level phenomena that cannot be observed directly, such as thin film growth and ion subplantation, and to examine the physical properties of nanotechnological devices that have not or cannot yet be created. In biophysics and structural biology, the method is frequently applied to study the motions of macromolecules such as proteins and nucleic acids, which can be useful for interpreting the results of certain biophysical experiments and for modeling interactions with other molecules, as in ligand docking. In principle, MD can be used for ab initio prediction of protein structure by simulating folding of the polypeptide chain from a random coil.
The results of MD simulations can be tested through comparison to experiments that measure molecular dynamics, of which a popular method is NMR spectroscopy. MD-derived structure predictions can be tested through community-wide experiments in Critical Assessment of Protein Structure Prediction (CASP), although the method has historically had limited success in this area. Michael Levitt, who shared the Nobel Prize partly for the application of MD to proteins, wrote in 1999 that CASP participants usually did not use the method due to "... a central embarrassment of molecular mechanics, namely that energy minimization or molecular dynamics generally leads to a model that is less like the experimental structure". Improvements in computational resources permitting more and longer MD trajectories, combined with modern improvements in the quality of force field parameters, have yielded some improvements in both structure prediction and homology model refinement, without reaching the point of practical utility in these areas; many identify force field parameters as a key area for further development.
MD simulation has been reported for pharmacophore development and drug design. For example, Pinto et al. implemented MD simulations of Bcl-xL complexes to calculate average positions of critical amino acids involved in ligand binding. Carlson et al. implemented molecular dynamics simulations to identify compounds that complement a receptor while causing minimal disruption to the conformation and flexibility of the active site. Snapshots of the protein at constant time intervals during the simulation were overlaid to identify conserved binding regions (conserved in at least three out of eleven frames) for pharmacophore development. Spyrakis et al. relied on a workflow of MD simulations, fingerprints for ligands and proteins (FLAP) and linear discriminant analysis (LDA) to identify the best ligand-protein conformations to act as pharmacophore templates based on retrospective ROC analysis of the resulting pharmacophores. In an attempt to ameliorate structure-based drug discovery modeling, vis-à-vis the need for many modeled compounds, Hatmal et al. proposed a combination of MD simulation and ligand-receptor intermolecular contacts analysis to discern critical intermolecular contacts (binding interactions) from redundant ones in a single ligand–protein complex. Critical contacts can then be converted into pharmacophore models that can be used for virtual screening.
An important factor is intramolecular hydrogen bonds, which are not explicitly included in modern force fields, but described as Coulomb interactions of atomic point charges. This is a crude approximation because hydrogen bonds have a partially quantum mechanical and chemical nature. Furthermore, electrostatic interactions are usually calculated using the dielectric constant of a vacuum, even though the surrounding aqueous solution has a much higher dielectric constant. Thus, using the macroscopic dielectric constant at short interatomic distances is questionable. Finally, van der Waals interactions in MD are usually described by Lennard-Jones potentials based on the Fritz London theory that is only applicable in a vacuum. However, all types of van der Waals forces are ultimately of electrostatic origin and therefore depend on dielectric properties of the environment. The direct measurement of attraction forces between different materials (as Hamaker constant) shows that "the interaction between hydrocarbons across water is about 10% of that across vacuum". The environment-dependence of van der Waals forces is neglected in standard simulations, but can be included by developing polarizable force fields.
Design constraints
The design of a molecular dynamics simulation should account for the available computational power. Simulation size (n = number of particles), timestep, and total time duration must be selected so that the calculation can finish within a reasonable time period. However, the simulations should be long enough to be relevant to the time scales of the natural processes being studied. To make statistically valid conclusions from the simulations, the time span simulated should match the kinetics of the natural process. Otherwise, it is analogous to making conclusions about how a human walks when only looking at less than one footstep. Most scientific publications about the dynamics of proteins and DNA use data from simulations spanning nanoseconds (10−9 s) to microseconds (10−6 s). To obtain these simulations, several CPU-days to CPU-years are needed. Parallel algorithms allow the load to be distributed among CPUs; an example is the spatial or force decomposition algorithm.
During a classical MD simulation, the most CPU intensive task is the evaluation of the potential as a function of the particles' internal coordinates. Within that energy evaluation, the most expensive one is the non-bonded or non-covalent part. In big O notation, common molecular dynamics simulations scale by if all pair-wise electrostatic and van der Waals interactions must be accounted for explicitly. This computational cost can be reduced by employing electrostatics methods such as particle mesh Ewald summation ( ), particle-particle-particle mesh (P3M), or good spherical cutoff methods ( ).
Another factor that impacts total CPU time needed by a simulation is the size of the integration timestep. This is the time length between evaluations of the potential. The timestep must be chosen small enough to avoid discretization errors (i.e., smaller than the period related to fastest vibrational frequency in the system). Typical timesteps for classical MD are on the order of 1 femtosecond (10−15 s). This value may be extended by using algorithms such as the SHAKE constraint algorithm, which fix the vibrations of the fastest atoms (e.g., hydrogens) into place. Multiple time scale methods have also been developed, which allow extended times between updates of slower long-range forces.
For simulating molecules in a solvent, a choice should be made between an explicit and implicit solvent. Explicit solvent particles (such as the TIP3P, SPC/E and SPC-f water models) must be calculated expensively by the force field, while implicit solvents use a mean-field approach. Using an explicit solvent is computationally expensive, requiring inclusion of roughly ten times more particles in the simulation. But the granularity and viscosity of explicit solvent is essential to reproduce certain properties of the solute molecules. This is especially important to reproduce chemical kinetics.
In all kinds of molecular dynamics simulations, the simulation box size must be large enough to avoid boundary condition artifacts. Boundary conditions are often treated by choosing fixed values at the edges (which may cause artifacts), or by employing periodic boundary conditions in which one side of the simulation loops back to the opposite side, mimicking a bulk phase (which may cause artifacts too).
Microcanonical ensemble (NVE)
In the microcanonical ensemble, the system is isolated from changes in moles (N), volume (V), and energy (E). It corresponds to an adiabatic process with no heat exchange. A microcanonical molecular dynamics trajectory may be seen as an exchange of potential and kinetic energy, with total energy being conserved. For a system of N particles with coordinates and velocities , the following pair of first order differential equations may be written in Newton's notation as
The potential energy function of the system is a function of the particle coordinates . It is referred to simply as the potential in physics, or the force field in chemistry. The first equation comes from Newton's laws of motion; the force acting on each particle in the system can be calculated as the negative gradient of .
For every time step, each particle's position and velocity may be integrated with a symplectic integrator method such as Verlet integration. The time evolution of and is called a trajectory. Given the initial positions (e.g., from theoretical knowledge) and velocities (e.g., randomized Gaussian), we can calculate all future (or past) positions and velocities.
One frequent source of confusion is the meaning of temperature in MD. Commonly we have experience with macroscopic temperatures, which involve a huge number of particles, but temperature is a statistical quantity. If there is a large enough number of atoms, statistical temperature can be estimated from the instantaneous temperature, which is found by equating the kinetic energy of the system to nkBT/2, where n is the number of degrees of freedom of the system.
A temperature-related phenomenon arises due to the small number of atoms that are used in MD simulations. For example, consider simulating the growth of a copper film starting with a substrate containing 500 atoms and a deposition energy of 100 eV. In the real world, the 100 eV from the deposited atom would rapidly be transported through and shared among a large number of atoms ( or more) with no big change in temperature. When there are only 500 atoms, however, the substrate is almost immediately vaporized by the deposition. Something similar happens in biophysical simulations. The temperature of the system in NVE is naturally raised when macromolecules such as proteins undergo exothermic conformational changes and binding.
Canonical ensemble (NVT)
In the canonical ensemble, amount of substance (N), volume (V) and temperature (T) are conserved. It is also sometimes called constant temperature molecular dynamics (CTMD). In NVT, the energy of endothermic and exothermic processes is exchanged with a thermostat.
A variety of thermostat algorithms are available to add and remove energy from the boundaries of an MD simulation in a more or less realistic way, approximating the canonical ensemble. Popular methods to control temperature include velocity rescaling, the Nosé–Hoover thermostat, Nosé–Hoover chains, the Berendsen thermostat, the Andersen thermostat and Langevin dynamics. The Berendsen thermostat might introduce the flying ice cube effect, which leads to unphysical translations and rotations of the simulated system.
It is not trivial to obtain a canonical ensemble distribution of conformations and velocities using these algorithms. How this depends on system size, thermostat choice, thermostat parameters, time step and integrator is the subject of many articles in the field.
Isothermal–isobaric (NPT) ensemble
In the isothermal–isobaric ensemble, amount of substance (N), pressure (P) and temperature (T) are conserved. In addition to a thermostat, a barostat is needed. It corresponds most closely to laboratory conditions with a flask open to ambient temperature and pressure.
In the simulation of biological membranes, isotropic pressure control is not appropriate. For lipid bilayers, pressure control occurs under constant membrane area (NPAT) or constant surface tension "gamma" (NPγT).
Generalized ensembles
The replica exchange method is a generalized ensemble. It was originally created to deal with the slow dynamics of disordered spin systems. It is also called parallel tempering. The replica exchange MD (REMD) formulation tries to overcome the multiple-minima problem by exchanging the temperature of non-interacting replicas of the system running at several temperatures.
Potentials in MD simulations
A molecular dynamics simulation requires the definition of a potential function, or a description of the terms by which the particles in the simulation will interact. In chemistry and biology this is usually referred to as a force field and in materials physics as an interatomic potential. Potentials may be defined at many levels of physical accuracy; those most commonly used in chemistry are based on molecular mechanics and embody a classical mechanics treatment of particle-particle interactions that can reproduce structural and conformational changes but usually cannot reproduce chemical reactions.
The reduction from a fully quantum description to a classical potential entails two main approximations. The first one is the Born–Oppenheimer approximation, which states that the dynamics of electrons are so fast that they can be considered to react instantaneously to the motion of their nuclei. As a consequence, they may be treated separately. The second one treats the nuclei, which are much heavier than electrons, as point particles that follow classical Newtonian dynamics. In classical molecular dynamics, the effect of the electrons is approximated as one potential energy surface, usually representing the ground state.
When finer levels of detail are needed, potentials based on quantum mechanics are used; some methods attempt to create hybrid classical/quantum potentials where the bulk of the system is treated classically but a small region is treated as a quantum system, usually undergoing a chemical transformation.
Empirical potentials
Empirical potentials used in chemistry are frequently called force fields, while those used in materials physics are called interatomic potentials.
Most force fields in chemistry are empirical and consist of a summation of bonded forces associated with chemical bonds, bond angles, and bond dihedrals, and non-bonded forces associated with van der Waals forces and electrostatic charge. Empirical potentials represent quantum-mechanical effects in a limited way through ad hoc functional approximations. These potentials contain free parameters such as atomic charge, van der Waals parameters reflecting estimates of atomic radius, and equilibrium bond length, angle, and dihedral; these are obtained by fitting against detailed electronic calculations (quantum chemical simulations) or experimental physical properties such as elastic constants, lattice parameters and spectroscopic measurements.
Because of the non-local nature of non-bonded interactions, they involve at least weak interactions between all particles in the system. Its calculation is normally the bottleneck in the speed of MD simulations. To lower the computational cost, force fields employ numerical approximations such as shifted cutoff radii, reaction field algorithms, particle mesh Ewald summation, or the newer particle–particle-particle–mesh (P3M).
Chemistry force fields commonly employ preset bonding arrangements (an exception being ab initio dynamics), and thus are unable to model the process of chemical bond breaking and reactions explicitly. On the other hand, many of the potentials used in physics, such as those based on the bond order formalism can describe several different coordinations of a system and bond breaking. Examples of such potentials include the Brenner potential for hydrocarbons and its
further developments for the C-Si-H and C-O-H systems. The ReaxFF potential can be considered a fully reactive hybrid between bond order potentials and chemistry force fields.
Pair potentials versus many-body potentials
The potential functions representing the non-bonded energy are formulated as a sum over interactions between the particles of the system. The simplest choice, employed in many popular force fields, is the "pair potential", in which the total potential energy can be calculated from the sum of energy contributions between pairs of atoms. Therefore, these force fields are also called "additive force fields". An example of such a pair potential is the non-bonded Lennard-Jones potential (also termed the 6–12 potential), used for calculating van der Waals forces.
Another example is the Born (ionic) model of the ionic lattice. The first term in the next equation is Coulomb's law for a pair of ions, the second term is the short-range repulsion explained by Pauli's exclusion principle and the final term is the dispersion interaction term. Usually, a simulation only includes the dipolar term, although sometimes the quadrupolar term is also included. When nl = 6, this potential is also called the Coulomb–Buckingham potential.
In many-body potentials, the potential energy includes the effects of three or more particles interacting with each other. In simulations with pairwise potentials, global interactions in the system also exist, but they occur only through pairwise terms. In many-body potentials, the potential energy cannot be found by a sum over pairs of atoms, as these interactions are calculated explicitly as a combination of higher-order terms. In the statistical view, the dependency between the variables cannot in general be expressed using only pairwise products of the degrees of freedom. For example, the Tersoff potential, which was originally used to simulate carbon, silicon, and germanium, and has since been used for a wide range of other materials, involves a sum over groups of three atoms, with the angles between the atoms being an important factor in the potential. Other examples are the embedded-atom method (EAM), the EDIP, and the Tight-Binding Second Moment Approximation (TBSMA) potentials, where the electron density of states in the region of an atom is calculated from a sum of contributions from surrounding atoms, and the potential energy contribution is then a function of this sum.
Semi-empirical potentials
Semi-empirical potentials make use of the matrix representation from quantum mechanics. However, the values of the matrix elements are found through empirical formulae that estimate the degree of overlap of specific atomic orbitals. The matrix is then diagonalized to determine the occupancy of the different atomic orbitals, and empirical formulae are used once again to determine the energy contributions of the orbitals.
There are a wide variety of semi-empirical potentials, termed tight-binding potentials, which vary according to the atoms being modeled.
Polarizable potentials
Most classical force fields implicitly include the effect of polarizability, e.g., by scaling up the partial charges obtained from quantum chemical calculations. These partial charges are stationary with respect to the mass of the atom. But molecular dynamics simulations can explicitly model polarizability with the introduction of induced dipoles through different methods, such as Drude particles or fluctuating charges. This allows for a dynamic redistribution of charge between atoms which responds to the local chemical environment.
For many years, polarizable MD simulations have been touted as the next generation. For homogenous liquids such as water, increased accuracy has been achieved through the inclusion of polarizability. Some promising results have also been achieved for proteins. However, it is still uncertain how to best approximate polarizability in a simulation. The point becomes more important when a particle experiences different environments during its simulation trajectory, e.g. translocation of a drug through a cell membrane.
Potentials in ab initio methods
In classical molecular dynamics, one potential energy surface (usually the ground state) is represented in the force field. This is a consequence of the Born–Oppenheimer approximation. In excited states, chemical reactions or when a more accurate representation is needed, electronic behavior can be obtained from first principles using a quantum mechanical method, such as density functional theory. This is named Ab Initio Molecular Dynamics (AIMD). Due to the cost of treating the electronic degrees of freedom, the computational burden of these simulations is far higher than classical molecular dynamics. For this reason, AIMD is typically limited to smaller systems and shorter times.
Ab initio quantum mechanical and chemical methods may be used to calculate the potential energy of a system on the fly, as needed for conformations in a trajectory. This calculation is usually made in the close neighborhood of the reaction coordinate. Although various approximations may be used, these are based on theoretical considerations, not on empirical fitting. Ab initio calculations produce a vast amount of information that is not available from empirical methods, such as density of electronic states or other electronic properties. A significant advantage of using ab initio methods is the ability to study reactions that involve breaking or formation of covalent bonds, which correspond to multiple electronic states. Moreover, ab initio methods also allow recovering effects beyond the Born–Oppenheimer approximation using approaches like mixed quantum-classical dynamics.
Hybrid QM/MM
QM (quantum-mechanical) methods are very powerful. However, they are computationally expensive, while the MM (classical or molecular mechanics) methods are fast but suffer from several limits (require extensive parameterization; energy estimates obtained are not very accurate; cannot be used to simulate reactions where covalent bonds are broken/formed; and are limited in their abilities for providing accurate details regarding the chemical environment). A new class of method has emerged that combines the good points of QM (accuracy) and MM (speed) calculations. These methods are termed mixed or hybrid quantum-mechanical and molecular mechanics methods (hybrid QM/MM).
The most important advantage of hybrid QM/MM method is the speed. The cost of doing classical molecular dynamics (MM) in the most straightforward case scales O(n2), where n is the number of atoms in the system. This is mainly due to electrostatic interactions term (every particle interacts with every other particle). However, use of cutoff radius, periodic pair-list updates and more recently the variations of the particle-mesh Ewald's (PME) method has reduced this to between O(n) to O(n2). In other words, if a system with twice as many atoms is simulated then it would take between two and four times as much computing power. On the other hand, the simplest ab initio calculations typically scale O(n3) or worse (restricted Hartree–Fock calculations have been suggested to scale ~O(n2.7)). To overcome the limit, a small part of the system is treated quantum-mechanically (typically active-site of an enzyme) and the remaining system is treated classically.
In more sophisticated implementations, QM/MM methods exist to treat both light nuclei susceptible to quantum effects (such as hydrogens) and electronic states. This allows generating hydrogen wave-functions (similar to electronic wave-functions). This methodology has been useful in investigating phenomena such as hydrogen tunneling. One example where QM/MM methods have provided new discoveries is the calculation of hydride transfer in the enzyme liver alcohol dehydrogenase. In this case, quantum tunneling is important for the hydrogen, as it determines the reaction rate.
Coarse-graining and reduced representations
At the other end of the detail scale are coarse-grained and lattice models. Instead of explicitly representing every atom of the system, one uses "pseudo-atoms" to represent groups of atoms. MD simulations on very large systems may require such large computer resources that they cannot easily be studied by traditional all-atom methods. Similarly, simulations of processes on long timescales (beyond about 1 microsecond) are prohibitively expensive, because they require so many time steps. In these cases, one can sometimes tackle the problem by using reduced representations, which are also called coarse-grained models.
Examples for coarse graining (CG) methods are discontinuous molecular dynamics (CG-DMD) and Go-models. Coarse-graining is done sometimes taking larger pseudo-atoms. Such united atom approximations have been used in MD simulations of biological membranes. Implementation of such approach on systems where electrical properties are of interest can be challenging owing to the difficulty of using a proper charge distribution on the pseudo-atoms. The aliphatic tails of lipids are represented by a few pseudo-atoms by gathering 2 to 4 methylene groups into each pseudo-atom.
The parameterization of these very coarse-grained models must be done empirically, by matching the behavior of the model to appropriate experimental data or all-atom simulations. Ideally, these parameters should account for both enthalpic and entropic contributions to free energy in an implicit way. When coarse-graining is done at higher levels, the accuracy of the dynamic description may be less reliable. But very coarse-grained models have been used successfully to examine a wide range of questions in structural biology, liquid crystal organization, and polymer glasses.
Examples of applications of coarse-graining:
protein folding and protein structure prediction studies are often carried out using one, or a few, pseudo-atoms per amino acid;
liquid crystal phase transitions have been examined in confined geometries and/or during flow using the Gay-Berne potential, which describes anisotropic species;
Polymer glasses during deformation have been studied using simple harmonic or FENE springs to connect spheres described by the Lennard-Jones potential;
DNA supercoiling has been investigated using 1–3 pseudo-atoms per basepair, and at even lower resolution;
Packaging of double-helical DNA into bacteriophage has been investigated with models where one pseudo-atom represents one turn (about 10 basepairs) of the double helix;
RNA structure in the ribosome and other large systems has been modeled with one pseudo-atom per nucleotide.
The simplest form of coarse-graining is the united atom (sometimes called extended atom) and was used in most early MD simulations of proteins, lipids, and nucleic acids. For example, instead of treating all four atoms of a CH3 methyl group explicitly (or all three atoms of CH2 methylene group), one represents the whole group with one pseudo-atom. It must, of course, be properly parameterized so that its van der Waals interactions with other groups have the proper distance-dependence. Similar considerations apply to the bonds, angles, and torsions in which the pseudo-atom participates. In this kind of united atom representation, one typically eliminates all explicit hydrogen atoms except those that have the capability to participate in hydrogen bonds (polar hydrogens). An example of this is the CHARMM 19 force-field.
The polar hydrogens are usually retained in the model, because proper treatment of hydrogen bonds requires a reasonably accurate description of the directionality and the electrostatic interactions between the donor and acceptor groups. A hydroxyl group, for example, can be both a hydrogen bond donor, and a hydrogen bond acceptor, and it would be impossible to treat this with one OH pseudo-atom. About half the atoms in a protein or nucleic acid are non-polar hydrogens, so the use of united atoms can provide a substantial savings in computer time.
Machine Learning Force Fields
Machine Learning Force Fields] (MLFFs) represent one approach to modeling interatomic interactions in molecular dynamics simulations. MLFFs can achieve accuracy close to that of ab initio methods. Once trained, MLFFs are much faster than direct quantum mechanical calculations. MLFFs address the limitations of traditional force fields by learning complex potential energy surfaces directly from high-level quantum mechanical data. Several software packages now support MLFFs, including VASP and open-source libraries like DeePMD-kit and SchNetPack.
Incorporating solvent effects
In many simulations of a solute-solvent system the main focus is on the behavior of the solute with little interest of the solvent behavior particularly in those solvent molecules residing in regions far from the solute molecule. Solvents may influence the dynamic behavior of solutes via random collisions and by imposing a frictional drag on the motion of the solute through the solvent. The use of non-rectangular periodic boundary conditions, stochastic boundaries and solvent shells can all help reduce the number of solvent molecules required and enable a larger proportion of the computing time to be spent instead on simulating the solute. It is also possible to incorporate the effects of a solvent without needing any explicit solvent molecules present. One example of this approach is to use a potential mean force (PMF) which describes how the free energy changes as a particular coordinate is varied. The free energy change described by PMF contains the averaged effects of the solvent.
Without incorporating the effects of solvent simulations of macromolecules (such as proteins) may yield unrealistic behavior and even small molecules may adopt more compact conformations due to favourable van der Waals forces and electrostatic interactions which would be dampened in the presence of a solvent.
Long-range forces
A long range interaction is an interaction in which the spatial interaction falls off no faster than where is the dimensionality of the system. Examples include charge-charge interactions between ions and dipole-dipole interactions between molecules. Modelling these forces presents quite a challenge as they are significant over a distance which may be larger than half the box length with simulations of many thousands of particles. Though one solution would be to significantly increase the size of the box length, this brute force approach is less than ideal as the simulation would become computationally very expensive. Spherically truncating the potential is also out of the question as unrealistic behaviour may be observed when the distance is close to the cut off distance.
Steered molecular dynamics (SMD)
Steered molecular dynamics (SMD) simulations, or force probe simulations, apply forces to a protein in order to manipulate its structure by pulling it along desired degrees of freedom. These experiments can be used to reveal structural changes in a protein at the atomic level. SMD is often used to simulate events such as mechanical unfolding or stretching.
There are two typical protocols of SMD: one in which pulling velocity is held constant, and one in which applied force is constant. Typically, part of the studied system (e.g., an atom in a protein) is restrained by a harmonic potential. Forces are then applied to specific atoms at either a constant velocity or a constant force. Umbrella sampling is used to move the system along the desired reaction coordinate by varying, for example, the forces, distances, and angles manipulated in the simulation. Through umbrella sampling, all of the system's configurations—both high-energy and low-energy—are adequately sampled. Then, each configuration's change in free energy can be calculated as the potential of mean force. A popular method of computing PMF is through the weighted histogram analysis method (WHAM), which analyzes a series of umbrella sampling simulations.
A lot of important applications of SMD are in the field of drug discovery and biomolecular sciences. For e.g. SMD was used to investigate the stability of Alzheimer's protofibrils, to study the protein ligand interaction in cyclin-dependent kinase 5 and even to show the effect of electric field on thrombin (protein) and aptamer (nucleotide) complex among many other interesting studies.
Examples of applications
Molecular dynamics is used in many fields of science.
First MD simulation of a simplified biological folding process was published in 1975. Its simulation published in Nature paved the way for the vast area of modern computational protein-folding.
First MD simulation of a biological process was published in 1976. Its simulation published in Nature paved the way for understanding protein motion as essential in function and not just accessory.
MD is the standard method to treat collision cascades in the heat spike regime, i.e., the effects that energetic neutron and ion irradiation have on solids and solid surfaces.
The following biophysical examples illustrate notable efforts to produce simulations of a systems of very large size (a complete virus) or very long simulation times (up to 1.112 milliseconds):
MD simulation of the full satellite tobacco mosaic virus (STMV) (2006, Size: 1 million atoms, Simulation time: 50 ns, program: NAMD) This virus is a small, icosahedral plant virus that worsens the symptoms of infection by Tobacco Mosaic Virus (TMV). Molecular dynamics simulations were used to probe the mechanisms of viral assembly. The entire STMV particle consists of 60 identical copies of one protein that make up the viral capsid (coating), and a 1063 nucleotide single stranded RNA genome. One key finding is that the capsid is very unstable when there is no RNA inside. The simulation would take one 2006 desktop computer around 35 years to complete. It was thus done in many processors in parallel with continuous communication between them.
Folding simulations of the Villin Headpiece in all-atom detail (2006, Size: 20,000 atoms; Simulation time: 500 μs= 500,000 ns, Program: Folding@home) This simulation was run in 200,000 CPU's of participating personal computers around the world. These computers had the Folding@home program installed, a large-scale distributed computing effort coordinated by Vijay Pande at Stanford University. The kinetic properties of the Villin Headpiece protein were probed by using many independent, short trajectories run by CPU's without continuous real-time communication. One method employed was the Pfold value analysis, which measures the probability of folding before unfolding of a specific starting conformation. Pfold gives information about transition state structures and an ordering of conformations along the folding pathway. Each trajectory in a Pfold calculation can be relatively short, but many independent trajectories are needed.
Long continuous-trajectory simulations have been performed on Anton, a massively parallel supercomputer designed and built around custom application-specific integrated circuits (ASICs) and interconnects by D. E. Shaw Research. The longest published result of a simulation performed using Anton is a 1.112-millisecond simulation of NTL9 at 355 K; a second, independent 1.073-millisecond simulation of this configuration was also performed (and many other simulations of over 250 μs continuous chemical time). In How Fast-Folding Proteins Fold, researchers Kresten Lindorff-Larsen, Stefano Piana, Ron O. Dror, and David E. Shaw discuss "the results of atomic-level molecular dynamics simulations, over periods ranging between 100 μs and 1 ms, that reveal a set of common principles underlying the folding of 12 structurally diverse proteins." Examination of these diverse long trajectories, enabled by specialized, custom hardware, allow them to conclude that "In most cases, folding follows a single dominant route in which elements of the native structure appear in an order highly correlated with their propensity to form in the unfolded state." In a separate study, Anton was used to conduct a 1.013-millisecond simulation of the native-state dynamics of bovine pancreatic trypsin inhibitor (BPTI) at 300 K.
Another important application of MD method benefits from its ability of 3-dimensional characterization and analysis of microstructural evolution at atomic scale.
MD simulations are used in characterization of grain size evolution, for example, when describing wear and friction of nanocrystalline Al and Al(Zr) materials. Dislocations evolution and grain size evolution are analyzed during the friction process in this simulation. Since MD method provided the full information of the microstructure, the grain size evolution was calculated in 3D using the Polyhedral Template Matching, Grain Segmentation, and Graph clustering methods. In such simulation, MD method provided an accurate measurement of grain size. Making use of these information, the actual grain structures were extracted, measured, and presented. Compared to the traditional method of using SEM with a single 2-dimensional slice of the material, MD provides a 3-dimensional and accurate way to characterize the microstructural evolution at atomic scale.
Molecular dynamics algorithms
Screened Coulomb potentials implicit solvent model
Integrators
Symplectic integrator
Verlet–Stoermer integration
Runge–Kutta integration
Beeman's algorithm
Constraint algorithms (for constrained systems)
Short-range interaction algorithms
Cell lists
Verlet list
Bonded interactions
Long-range interaction algorithms
Ewald summation
Particle mesh Ewald summation (PME)
Particle–particle-particle–mesh (P3M)
Shifted force method
Parallelization strategies
Domain decomposition method (Distribution of system data for parallel computing)
Ab-initio molecular dynamics
Car–Parrinello molecular dynamics
Specialized hardware for MD simulations
Anton – A specialized, massively parallel supercomputer designed to execute MD simulations
MDGRAPE – A special purpose system built for molecular dynamics simulations, especially protein structure prediction
Graphics card as a hardware for MD simulations
| Physical sciences | Molecular physics | Physics |
12254052 | https://en.wikipedia.org/wiki/Smoking | Smoking | Smoking is a practice in which a substance is combusted and the resulting smoke is typically inhaled to be tasted and absorbed into the bloodstream of a person. Most commonly, the substance used is the dried leaves of the tobacco plant, which have been rolled with a small rectangle of paper into an elongated cylinder called a cigarette. Other forms of smoking include the use of a smoking pipe or a bong.
Smoking is primarily practised as a route of administration for psychoactive chemicals because the active substances within the burnt dried plant leaves vaporize and can be airborne-delivered into the respiratory tract, where they are rapidly absorbed into the bloodstream of the lungs and then reach the central nervous system. In the case of tobacco smoking, these active substances are a mixture of aerosol particles that includes the pharmacologically active alkaloid nicotine, which stimulates the nicotinic acetylcholine receptors in the brain. Other notable active substances inhaled via smoking include tetrahydrocannabinol (from cannabis), morphine (from opium) and cocaine (from crack).
Smoking is one of the most common forms of recreational drug use. Tobacco smoking is the most popular form, being practised by over one billion people globally, of whom the majority are in the developing countries. Less common drugs for smoking include cannabis and opium. Some of the substances are classified as hard narcotics, like heroin, but the use of these is very limited as they are usually not commercially available. Cigarettes are primarily industrially manufactured but also can be hand-rolled from loose tobacco and rolling paper. Other smoking implements include pipes, cigars, bidis, hookahs, and bongs.
Smoking has negative health effects, because smoke inhalation inherently poses challenges to various physiologic processes such as respiration. Smoking tobacco is among the leading causes of many diseases such as lung cancer, heart attack, COPD, erectile dysfunction, and birth defects. Diseases related to tobacco smoking have been shown to kill approximately half of long-term smokers when compared to average mortality rates faced by non-smokers. Smoking caused over five million deaths a year from 1990 to 2015. Non-smokers account for 600,000 deaths globally due to second-hand smoke. The health hazards of smoking have caused many countries to institute high taxes on tobacco products, publish advertisements to discourage use, limit advertisements that promote use, and provide help with quitting for those who do smoke.
Smoking can be dated to as early as 5000 BCE, and has been recorded in many different cultures across the world. Early smoking evolved in association with religious ceremonies; as offerings to deities; in cleansing rituals; or to allow shamans and priests to alter their minds for purposes of divination or spiritual enlightenment. After the European exploration and conquest of the Americas, the practice of smoking tobacco quickly spread to the rest of the world. In regions like India and Sub-Saharan Africa, it merged with existing practices of smoking (mostly of cannabis). In Europe, it introduced a new type of social activity and a form of drug intake which previously had been unknown.
Perception surrounding smoking has varied over time and from one place to another: holy and sinful, sophisticated and vulgar, a panacea and deadly health hazard. In the last decade of the 20th century, smoking came to be viewed in a decidedly negative light, especially in Western countries.
History
Early uses
The history of smoking dates back to as early as 5000 BCE for shamanistic rituals. Many ancient civilizations, such as the Babylonian and Chinese, burnt incense as a part of religious rituals, as did the Israelites and the later Catholic and Orthodox Christian churches. Smoking in the Americas probably had its origins in the incense-burning ceremonies of shamans but was later adopted for pleasure, or as a social tool. The smoking of tobacco, as well as various hallucinogenic drugs, was used to achieve trances and to come into contact with the spirit world.
Substances such as cannabis, clarified butter (ghee), fish offal, dried snake skins and various pastes molded around incense sticks dates back at least 2000 years. Fumigation (dhupa) and fire offerings (homa) are prescribed in the Ayurveda for medical purposes, and have been practiced for at least 3,000 years while smoking, dhumrapana (literally "drinking smoke"), has been practiced for at least 2,000 years. Before modern times these substances have been consumed through pipes, with stems of various lengths or chillums. Archaeological findings also show the existence of pipes to smoke opium in Cyprus and Crete as soon as the Bronze Age.
Cannabis smoking was common in the Middle East before the arrival of tobacco, and was early on a common social activity that centered around the type of water pipe called a hookah. Smoking, especially after the introduction of tobacco, was an essential component of Muslim society and culture and became integrated with important traditions such as weddings, funerals and was expressed in architecture, clothing, literature and poetry.
Cannabis smoking was introduced to Sub-Saharan Africa through Ethiopia and the east African coast by either Indian or Arab traders in the 13th century or earlier and spread on the same trade routes as those that carried coffee, which originated in the highlands of Ethiopia. It was smoked in calabash water pipes with terracotta smoking bowls, apparently an Ethiopian invention which was later conveyed to eastern, southern and central Africa.
Reports from the first European explorers and conquistadors to reach the Americas tell of rituals where native priests smoked themselves into such high degrees of intoxication that it is unlikely that the rituals were limited to just tobacco.
Popularization
In 1612, six years after the settlement of Jamestown, John Rolfe was credited as the first settler to successfully grow tobacco as a cash crop. The demand quickly grew as tobacco, referred to as "golden weed", revived the Virginia Company from its failed expeditions in search for gold in the Americas. In order to meet demands from the old world, tobacco was grown in succession, quickly depleting the land. This became a motivator to settle west into the unknown continent, and likewise an expansion of tobacco production. Indentured servants became the primary labor force up until Bacon's Rebellion, from which the focus turned to slavery. This trend abated following the American Revolution as slavery became regarded as unprofitable. However the practice was revived in 1794 with the invention of the cotton gin.
A Frenchman named Jean Nicot (from whose name the word nicotine is derived) introduced tobacco to France in 1560. From France tobacco spread to England. The first report documents an English sailor in Bristol in 1556, seen "emitting smoke from his nostrils". Like tea, coffee and opium, tobacco was just one of many intoxicants that was originally used as a form of medicine. Tobacco was introduced around 1600 by French merchants in what today is modern-day The Gambia and Senegal. At the same time caravans from Morocco brought tobacco to the areas around Timbuktu and the Portuguese brought the commodity (and the plant) to southern Africa, establishing the popularity of tobacco throughout all of Africa by the 1650s.
Soon after its introduction to the Old World, tobacco came under frequent criticism from state and religious leaders. Murad IV, sultan of the Ottoman Empire 1623–40 was among the first to attempt a smoking ban by claiming it was a threat to public morality and health. The Chongzhen Emperor of China issued an edict banning smoking two years before his death and the overthrow of the Ming dynasty. Later, the Manchu rulers of the Qing dynasty, would proclaim smoking "a more heinous crime than that even of neglecting archery". In Edo period Japan, some of the earliest tobacco plantations were scorned by the shōgun as being a threat to the military economy by letting valuable farmland go to waste for the use of a recreational drug instead of being used to plant food crops.
Religious leaders have often been prominent among those who considered smoking immoral or outright blasphemous. In 1634, the Patriarch of Moscow and all Rus' forbade the sale of tobacco and sentenced men and women who flouted the ban to have their nostrils slit and their backs whipped until skin came off their backs. The Western church leader Pope Urban VII likewise condemned smoking in a papal bull of 1590. Despite many concerted efforts, restrictions and bans were almost universally ignored. When James VI and I, a staunch anti-smoker and the author of A Counterblaste to Tobacco, tried to curb the new trend by enforcing a whopping 4000% tax increase on tobacco in 1604, it proved a failure, as London had some 7,000 tobacco sellers by the early 17th century. Later, scrupulous rulers would realise the futility of smoking bans and instead turned tobacco trade and cultivation into lucrative government monopolies.
By the mid-17th century every major civilization had been introduced to tobacco smoking and in many cases had already assimilated it into its culture, despite the attempts of many rulers to stamp the practice out with harsh penalties or fines. Tobacco, both product, and plant followed the major trade routes to major ports and markets, and then on into the hinterlands. The English language term smoking was coined in the late 18th century; before then the practice was referred to as drinking smoke.
Tobacco and cannabis were used in Sub-Saharan Africa, much like elsewhere in the world, to confirm social relations, but also created entirely new ones. In what is today Congo, a society called Bena Diemba ("People of Cannabis") was organized in the late 19th century in Lubuko ("The Land of Friendship"). The Bena Diemba were collectivist pacifists that rejected alcohol and herbal medicines in favor of cannabis.
The growth remained stable until the American Civil War in the 1860s, from which the primary labor force transition from slavery to sharecropping. This compounded with a change in demand, lead to the industrialization of tobacco production with the cigarette. James Albert Bonsack, a craftsman, in 1881 produced a machine to speed the production of cigarettes.
Opium
In the 19th century, the practise of smoking opium became widespread in China. Previously, opium had only been ingested via consumption, and then only for its medicinal properties (opium was an anaesthetic). The narcotic was also outlawed in China sometime in the early 18th century due the societal issues it caused. Due to a massive trade imbalance, however, foreign merchants started to smuggle opium into China via Canton, to the chagrin of the Chinese authorities. Attempts by Chinese official Lin Zexu to eliminate the trade led to the outbreak of the First Opium War. The Chinese defeat in the First and Second Opium Wars resulted in the legalization of the importation of opium into China.
Opium smoking later spread with Chinese immigrants and spawned many infamous opium dens in Chinatowns around South and Southeast Asia, Europe and the Americas. In the latter half of the 19th century, opium smoking became popular in the artistic community in Europe, especially Paris; artists' neighborhoods such as Montparnasse and Montmartre became virtual "opium capitals". While opium dens that catered primarily to emigrant Chinese continued to exist in Chinatowns around the world, the trend among the European artists largely abated after the outbreak of World War I. The consumption of Opium abated in China during the Cultural Revolution in the 1960s and 1970s.
Anti-tobacco movement
Many people have been critical about tobacco use since it gained popularity. In 1798, Dr. Benjamin Rush (early American physician, signer of the Declaration of Independence, Surgeon General under George Washington, and anti-tobacco activist) was "against the habitual use of tobacco" because he believed it (a) "led to a desire for strong drink," (b) "was injurious both to health and morals," (c) "is generally offensive to" nonsmokers, (d) "produces a want of respect for" nonsmokers, and (e) "always disposes to unkind and unjust behavior towards them."
With the modernization of cigarette production compounded with the increased life expectancies during the 1920s, adverse health effects began to become more prevalent. In Germany, anti-smoking groups, often associated with anti-liquor groups, first published advocacy against the consumption of tobacco in the journal Der Tabakgegner (The Tobacco Opponent) in 1912 and 1932. In 1929, Fritz Lickint of Dresden, Germany, published a paper containing formal statistical evidence of a lung cancer–tobacco link. During the Great Depression, Adolf Hitler condemned his earlier smoking habit as a waste of money, and later with stronger assertions. This movement was further strengthened with Nazi reproductive policy as women who smoked were viewed as unsuitable to be wives and mothers in a German family.
The movement in Nazi Germany did reach across enemy lines during the Second World War, as anti-smoking groups quickly lost popular support. By the end of the Second World War, American cigarette manufacturers quickly reentered the German black market. Illegal smuggling of tobacco became prevalent, and leaders of the Nazi anti-smoking campaign were assassinated. As part of the Marshall Plan, the United States shipped free tobacco to Germany; with 24,000 tons in 1948 and 69,000 tons in 1949. Per capita yearly cigarette consumption in post-war Germany steadily rose from 460 in 1950 to 1,523 in 1963. By the end of the 20th century, anti-smoking campaigns in Germany were unable to exceed the effectiveness of the Nazi-era climax in the years 1939–41 and German tobacco health research was described by Robert N. Proctor as "muted".
In the UK and the US, an increase in lung cancer rates, formerly "among the rarest forms of disease", was noted by the 1930s, but its cause remained unknown and even the credibility of this increase was sometimes disputed as late as 1950. For example, in Connecticut, reported age-adjusted incidence rates of lung cancer among males increased 220% between 1935–39 and 1950–54. In the UK, the share of lung cancer among all cancer deaths in men increased from 1.5% in 1920 to 19.7% in 1947. Nevertheless, these increases were questioned as potentially caused by increased reporting and improved methods of diagnosis. Although several carcinogens were already known at the time (for example, benzo[a]pyrene was isolated from coal tar and demonstrated to be a potent carcinogen in 1933), none were known to be contained in adequate quantities in tobacco smoke. Richard Doll in 1950 published research in the British Medical Journal showing a close link between smoking and lung cancer. Four years later, in 1954 the British Doctors Study, a study of some 40 thousand doctors over 20 years, confirmed the link, based on which the government issued advice that smoking and lung cancer rates were related. In 1964 the United States Surgeon General's Report on Smoking and Health demonstrated the relationship between smoking and cancer. Further reports confirmed this link in the 1980s and concluded in 1986 that passive smoking was also harmful.
As scientific evidence mounted in the 1980s, tobacco companies claimed contributory negligence as the adverse health effects were previously unknown or lacked substantial credibility. Health authorities sided with these claims up until 1998, from which they reversed their position. The Tobacco Master Settlement Agreement, originally between the four largest US tobacco companies and the Attorneys General of 46 states, restricted certain types of tobacco advertisement and required payments for health compensation; which later amounted to the largest civil settlement in United States history.
From 1965 to 2006, rates of smoking in the United States have declined from 42% to 20.8%. A significant majority of those who quit were professional, affluent men. Despite this decrease in the prevalence of consumption, the average number of cigarettes consumed per person per day increased from 22 in 1954 to 30 in 1978. This paradoxical event suggests that those who quit smoked less, while those who continued to smoke moved to smoke more light cigarettes. This trend has been paralleled by many industrialized nations as rates have either leveled-off or declined. In the developing countries, however, tobacco consumption continues to rise at 3.4% in 2002. In Africa, smoking is in most areas considered to be modern, and many of the strong adverse opinions that prevail in the West receive much less attention. Today Russia leads as the top consumer of tobacco followed by Indonesia, Laos, Ukraine, Belarus, Greece, Jordan, and China.
At the global scale, initial ideas of an international convention towards the prevention of tobacco had been initiated in the World Health Assembly (WHA) in 1996. In 1998, along with the successful election of Dr. Gro Harlem Brundtland as the Director-General, the World Health Organization set tobacco control as its leading health concern and has begun a program known as the Tobacco Free Initiative (TFI) in order to reduce rates of consumption in the developing world. However, it was not until 2003 that the Framework Convention on Tobacco Control (FCTC) was accepted in WHA and entered into force in 2005. FCTC marked a milestone as the first international treaty concerning a global health issue that aims to combat tobacco in multiple aspects including tobacco taxes, advertisement, trading, environmental affects, health influences, etc. The birth of this evidence-based and systematic approach has resulted in the reinforcement of tobacco taxes and the implementation of smoke-free laws in 128 countries that led to the decrease of smoking prevalence in developing nations. In Nepal, "Smokers are not selfish", a health campaign lasting two weeks is started on the occasion of Valentine day and Vasant panchami to motiviate individuals to quit smoking as a sacrifice for their loved ones and making it a meaningful decision of life. This campaign is attracting public attention.
Other substances
In the early 1980s, organized international trafficking of cocaine grew. However, overproduction and tighter legal enforcement for the illegal product caused drug dealers to convert the powder to "crack" – a solid, smokable form of cocaine that could be sold in smaller quantities to more people. This trend abated in the 1990s as increased police action coupled with a robust economy caused many potential consumers to give up or fail to take up the habit.
Recent years shows an increase in the consumption of vaporized heroin, methamphetamine and Phencyclidine (PCP). Along with a smaller number of psychedelic drugs such as Changa, DMT, 5-Meo-DMT, and Salvia divinorum.
Substances and equipment
The most popular type of substance that is smoked is tobacco. There are many different tobacco cultivars which are made into a wide variety of mixtures and brands. Tobacco is often sold flavored, often with various fruit aromas, something which is especially popular for use with water pipes, such as hookahs. The second most common substance that is smoked is cannabis, made from the flowers or leaves of Cannabis sativa or Cannabis indica. The substance is considered illegal in most countries in the world and in the countries that tolerate public consumption, it is sometimes only pseudo-legal. Despite this, a considerable percentage of the adult population in many countries have tried it with smaller minorities doing it on a regular basis. Since cannabis is illegal or only tolerated in many jurisdictions, there is no industrial mass-production of cigarettes, meaning that the most common form of smoking is with hand-rolled cigarettes (often called joints) or with pipes. Water pipes are also fairly common; water pipes used for cannabis include designs known as bongs and bubblers, among others.
A few other recreational drugs are smoked by smaller minorities. Most of these substances are controlled, and some are considerably more intoxicating than either tobacco or cannabis. These include crack cocaine, heroin, methamphetamine and PCP. A small number of psychedelic drugs are also smoked, including DMT, 5-Meo-DMT, and Salvia divinorum.
Even the most primitive form of smoking requires tools of some sort to perform. This has resulted in a staggering variety of smoking tools and paraphernalia from all over the world. Whether tobacco, cannabis, opium or herbs, some form of receptacle is required along with a source of fire to light the mixture. The most common today is by far the cigarette, consisting of a mild inhalant strain of tobacco in a tightly rolled tube of paper, usually manufactured industrially and including a filter, or hand-rolled with loose tobacco. Other popular smoking tools are various pipes and cigars.
A less common but increasingly popular alternative to smoking is vaporizers, which use hot air convection to deliver the substance without combustion, which may reduce health risks. A portable vaporization alternative appeared in 2003 with the introduction of electronic cigarettes, battery-operated, cigarette-shaped devices which produce an aerosol intended to mimic the smoke from burning tobacco, delivering nicotine to the user without some of the harmful substances released in tobacco smoke.
Other than actual smoking equipment, many other items are associated with smoking; cigarette cases, cigar boxes, lighters, matchboxes, cigarette holders, cigar holders, ashtrays, silent butlers, pipe cleaners, tobacco cutters, match stands, pipe tampers, cigarette companions and so on. Some examples of these have become valuable collector items and particularly ornate and antique items can fetch high prices.
Health effects
Smoking is one of the leading preventable causes of deaths globally and is the cause of over 8 million deaths annually, 1.2 million of which are non-smokers who die due to second-hand smoke. In the United States, about 500,000 deaths per year are attributed to smoking-related diseases and a recent study estimated that as much as one-third of China's male population will have significantly shortened lifespans due to smoking. Male and female smokers lose an average of 13.2 and 14.5 years of life, respectively. At least half of all lifelong smokers die earlier as a result of smoking. The risk of dying from lung cancer before age 85 is 22.1% for a male smoker and 11.9% for a female current smoker, in the absence of competing causes of death. The corresponding estimates for lifelong nonsmokers are a 1.1% probability of dying from lung cancer before age 85 for a man of European descent, and a 0.8% probability for a woman. Smoking just one cigarette a day results in a risk of coronary heart disease that is halfway between that of a heavy smoker and a non-smoker. The non-linear dose–response relationship may be explained by smoking's effect on platelet aggregation.
Among the diseases that can be caused by smoking are vascular stenosis, lung cancer, heart attacks and chronic obstructive pulmonary disease (COPD). Smoking during pregnancy may cause ADHD to a fetus.
Smoking is a risk factor strongly associated with periodontitis and tooth loss. The effects of smoking on periodontal tissues depend on the number of cigarettes smoked daily and the duration of the habit. A study showed that smokers had 2.7 times and former smokers 2.3 times greater probabilities to have established periodontal disease than non‐smokers, independent of age, sex and plaque index, however, the effect of tobacco on periodontal tissues seems to be more pronounced in men than in women. Studies have found that smokers had greater odds for more severe dental bone loss compared to non‐smokers; also, people who smoke and drink alcohol heavily have much higher risk of developing oral cancer (mouth and lip) compared with people who do neither. Smoking can also cause milanosis in the mouth.
Smoking has been also associated with oral conditions including dental caries, dental implant failures, premalignant lesions, and cancer. Smoking can affect the immune-inflammatory processes which may increase susceptibility to infections; it can alter the oral mycobiota and facilitate colonization of the oral cavity with fungi and pathogenic molds.
Many governments are trying to deter people from smoking with anti-smoking campaigns in mass media stressing the harmful long-term effects of smoking. Passive smoking, or secondhand smoking, which affects people in the immediate vicinity of smokers, is a major reason for the enforcement of smoking bans. These are laws enforced to stop individuals from smoking in indoor public places, such as bars, pubs and restaurants, thus reducing nonsmokers' exposure to secondhand smoke. A common concern among legislators is to discourage smoking among minors and many states have passed laws against selling tobacco products to underage customers (establishing a smoking age). Many developing countries have not adopted anti-smoking policies, leading some to call for anti-smoking campaigns and further education to explain the negative effects of ETS (Environmental Tobacco Smoke) in developing countries. Tobacco advertising is also sometimes regulated to make smoking less appealing.
Despite the many bans, European countries still hold 18 of the top 20 spots, and according to the ERC, a market research company, the heaviest smokers are from Greece, averaging 3,000 cigarettes per person in 2007. Rates of smoking have leveled off or declined in the developed world but continue to rise in developing countries. Smoking rates in the United States have dropped by half from 1965 to 2006, falling from 42% to 20.8% in adults.
The effects of addiction on society vary considerably between different substances that can be smoked and the indirect social problems that they cause, in great part because of the differences in legislation and the enforcement of narcotics legislation around the world. Though nicotine is a highly addictive drug, its effects on cognition are not as intense or noticeable as other drugs such as cocaine, amphetamines or any of the opiates (including heroin and morphine).
Smoking is a risk factor in Alzheimer's disease. While smoking more than 15 cigarettes per day has been shown to worsen the symptoms of Crohn's disease, smoking has been shown to actually lower the prevalence of ulcerative colitis.
Smokers are 30-40% more likely to develop type 2 diabetes than non-smokers, and the risk increases with the number of cigarettes smoked.
Physiology
Inhaling the vaporized gas form of substances into the lungs is a quick and very effective way of delivering drugs into the bloodstream (as the gas diffuses directly into the pulmonary vein, then into the heart and from there to the brain) and affects the user within less than a second of the first inhalation. The lungs consist of several million tiny bulbs called alveoli that altogether have an area of over 70 m2 (about the area of a tennis court). This can be used to administer useful medical as well as recreational drugs such as aerosols, consisting of tiny droplets of a medication, or as gas produced by burning plant material with a psychoactive substance or pure forms of the substance itself. Not all drugs can be smoked, for example the sulphate derivative that is most commonly inhaled through the nose, though purer free base forms of substances can, but often require considerable skill in administering the drug properly. The method is also somewhat inefficient since not all of the smoke will be inhaled. The inhaled substances trigger chemical reactions in nerve endings in the brain due to being similar to naturally occurring substances such as endorphins and dopamine, which are associated with sensations of pleasure. The result is what is usually referred to as a "high" that ranges between the mild stimulus caused by nicotine to the intense euphoria caused by heroin, cocaine and methamphetamines.
Inhaling smoke into the lungs, no matter the substance, has adverse effects on one's health. The incomplete combustion produced by burning plant material, like tobacco or cannabis, produces carbon monoxide, which impairs the ability of blood to carry oxygen when inhaled into the lungs. There are several other toxic compounds in tobacco that constitute serious health hazards to long-term smokers from a whole range of causes; vascular abnormalities such as stenosis, lung cancer, heart attacks, strokes, impotence, low birth weight of infants born by smoking mothers. 8% of long-term smokers develop the characteristic set of facial changes known to doctors as smoker's face.
Tobacco smoke is a complex mixture of over 5,000 identified chemicals, of which 98 are known to have specific toxicological properties. The most important chemicals causing cancer are those that produce DNA damage since such damage appears to be the primary underlying cause of cancer. Cunningham et al. combined the microgram weight of the compound in the smoke of one cigarette with the known genotoxic effect per microgram to identify the most carcinogenic compounds in cigarette smoke. The seven most important carcinogens in tobacco smoke are shown in the table, along with DNA alterations they cause.
Psychology
Most tobacco smokers begin during adolescence or early adulthood. Smoking has elements of risk-taking and rebellion, which often appeal to young people. The presence of high-status models and peers may also encourage smoking. Because teenagers are influenced more by their peers than by adults, attempts by parents, schools, and health professionals at preventing people from trying cigarettes are not always successful.
Smokers often report that cigarettes help relieve feelings of stress. However, the stress levels of adult smokers are slightly higher than those of nonsmokers. Adolescent smokers report increasing levels of stress as they develop regular patterns of smoking, and smoking cessation leads to reduced stress. Far from acting as an aid for mood control, nicotine dependency seems to exacerbate stress. This is confirmed in the daily mood patterns described by smokers, with normal moods during smoking and worsening moods between cigarettes. Thus, the apparent relaxant effect of smoking only reflects the reversal of the tension and irritability that develop during nicotine depletion. Dependent smokers need nicotine to remain feeling normal.
In the mid-20th century psychologists such as Hans Eysenck developed a personality profile for the typical smoker of that period; extraversion was associated with smoking, and smokers tended to be sociable, impulsive, risk taking, and excitement-seeking individuals. Although personality and social factors may make people likely to smoke, the actual habit is a function of operant conditioning. During the early stages, smoking provides pleasurable sensations (because of its action on the dopamine system) and thus serves as a source of positive reinforcement. After an individual has smoked for many years, the avoidance of withdrawal symptoms and negative reinforcement become the key motivations. Like all addictive substances, the amount of exposure required to become dependent on nicotine can vary from person to person.
In terms of the Big Five personality traits, research has found smoking to be correlated with lower levels of agreeableness and conscientiousness, as well as higher levels of extraversion and neuroticism.
Prevention
Education and counselling by physicians of children and adolescents has been found to be effective in decreasing the risk of tobacco use. Systematic reviews show that psychosocial interventions can help women stop smoking in late pregnancy, reducing low birthweight and preterm births. A 2016 Cochrane review showed that the combination of medication and behavioural support was more effective than minimal interventions or usual care. Another Cochrane review "suggests that neither reducing smoking to quit nor quitting abruptly results in superior quit rates; people could therefore be given a choice of how to quit, and support provided to people who would specifically like to reduce their smoking before quitting."
Prevalence
Smoking, primarily of tobacco, is an activity that is practiced by some 1.1 billion people, and up to 1/3 of the adult population. The image of the smoker can vary considerably, but is very often associated, especially in fiction, with individuality and aloofness. Even so, smoking of both tobacco and cannabis can be a social activity which serves as a reinforcement of social structures and is part of the cultural rituals of many and diverse social and ethnic groups. Many smokers begin smoking in social settings and the offering and sharing of a cigarette is often an important rite of initiation or simply a good excuse to start a conversation with strangers in many settings; in bars, night clubs, at work or on the street. Lighting a cigarette is often seen as an effective way of avoiding the appearance of idleness or mere loitering. For adolescents, it can function as a first step out of childhood or as an act of rebellion against the adult world. Also, smoking can be seen as a sort of camaraderie. It has been shown that even opening a packet of cigarettes, or offering a cigarette to other people, can increase the level of dopamine (the "happy feeling") in the brain, and it is doubtless that people who smoke form relationships with fellow smokers, in a way that only proliferates the habit, particularly in countries where smoking inside public places has been made illegal. Other than recreational drug use, it can be used to construct identity and a development of self-image by associating it with personal experiences connected with smoking. The rise of the modern anti-smoking movement in the late 19th century did more than create awareness of the hazards of smoking; it provoked reactions of smokers against what was, and often still is, perceived as an assault on personal freedom and has created an identity among smokers as rebels or outcasts, apart from non-smokers:
The importance of tobacco to soldiers was early on recognized as something that could not be ignored by commanders. By the 17th century allowances of tobacco were a standard part of the naval rations of many nations and by World War I cigarette manufacturers and governments collaborated in securing tobacco and cigarette allowances to soldiers in the field. It was asserted that regular use of tobacco while under duress would not only calm the soldiers but allow them to withstand greater hardship. Until the mid-20th century, the majority of the adult population in many Western nations were smokers and the claims of anti-smoking activists were met with much skepticism, if not outright contempt. Today the movement has considerably more weight and evidence of its claims, but a considerable proportion of the population remains steadfast smokers.
Society and culture
Smoking has been accepted into culture, in various art forms, and has developed many distinct, and often conflicting or mutually exclusive, meanings depending on time, place and the practitioners of smoking. Pipe smoking, until recently one of the most common forms of smoking, is today often associated with solemn contemplation, old age and is often considered quaint and archaic. Cigarette smoking, which did not begin to become widespread until the late 19th century, has more associations of modernity and the faster pace of the industrialized world. Cigars have been, and still are, associated with masculinity, power and is an iconic image associated with the stereotypical capitalist. In fact, some evidence suggests that men with higher than average testosterone levels are more likely to smoke. Smoking in public has for a long time been something reserved for men and when done by women has been associated with promiscuity. In Japan during the Edo period, prostitutes and their clients would often approach one another under the guise of offering a smoke; the same was true for 19th-century Europe.
Art
The earliest depictions of smoking can be found on Classical Mayan pottery from around the 9th century. The art was primarily religious in nature and depicted deities or rulers smoking early forms of cigarettes. Soon after smoking was introduced outside of the Americas it began appearing in painting in Europe and Asia. The painters of the Dutch Golden Age were among the first to paint portraits of people smoking and still lifes of pipes and tobacco. For southern European painters of the 17th century, a pipe was much too modern to include in the preferred motifs inspired by mythology from Greek and Roman antiquity. At first smoking was considered lowly and was associated with peasants. Many early paintings were of scenes set in taverns or brothels. Later, as the Dutch Republic rose to considerable power and wealth, smoking became more common amongst the affluent and portraits of elegant gentlemen tastefully raising a pipe appeared. Smoking represented pleasure, transience and the briefness of earthly life as it, quite literally, went up in smoke. Smoking was also associated with representations of both the sense of smell and that of taste.
In the 18th century smoking became far more sparse in painting as the elegant practice of taking snuff became popular. Smoking a pipe was again relegated to portraits of lowly commoners and country folk and the refined sniffing of shredded tobacco followed by sneezing was rare in art. When smoking appeared it was often in the exotic portraits influenced by Orientalism. Many proponents of postcolonialism controversially believe this portrayal was a means of projecting an image of European superiority over its colonies and a perception of the male dominance of a feminized Orient. Proponents believe the theme of the exotic and alien "Other" escalated in the 19th century, fueled by the rise in the popularity of ethnology during the Enlightenment.
In the 19th century smoking was common as a symbol of simple pleasures; the pipe smoking "noble savage", solemn contemplation by Classical Roman ruins, scenes of an artist becoming one with nature while slowly toking a pipe. The newly empowered middle class also found a new dimension of smoking as a harmless pleasure enjoyed in smoking saloons and libraries. Smoking a cigarette or a cigar would also become associated with the Bohemian, someone who shunned the conservative middle class values and displayed his contempt for conservatism. But this was a pleasure that was to be confined to a male world; women smokers were associated with prostitution and smoking was not considered an activity fit for proper ladies. It was not until the start of the 20th century that smoking women would appear in paintings and photos, giving a chic and charming impression. Impressionists like Vincent van Gogh, who was a pipe smoker himself, would also begin to associate smoking with gloom and fin-du-siècle fatalism.
While the symbolism of the cigarette, pipe and cigar respectively were consolidated in the late 19th century, it was not until the 20th century that artists began to use it fully; a pipe would stand for thoughtfulness and calm; the cigarette symbolized modernity, strength and youth, but also nervous anxiety; the cigar was a sign of authority, wealth and power. The decades following World War II, during the apex of smoking when the practice had still not come under fire by the growing anti-smoking movement, a cigarette casually tucked between the lips represented the young rebel, epitomized in actors like Marlon Brando and James Dean or mainstays of advertising like the Marlboro Man. It was not until the 1970s when the negative aspects of smoking began to appear, yielding the image of the unhealthy lower-class individual, reeking of cigarette smoke and lack of motivation and drive, which was especially prominent in art inspired or commissioned by anti-smoking campaigns. In his painting "Holy Smokes", artist Brian Whelan pokes fun at the smoking debate and its newly found focus on morality and guilt.
Film and TV
Ever since the era of silent films, smoking has had a major part in film symbolism. In the hard-boiled film noir crime thrillers, cigarette smoke often frames characters and is frequently used to add an aura of mystique or nihilism. One of the forerunners of this symbolism can be seen in Fritz Lang's Weimar era Dr Mabuse, der Spieler, 1922 (Dr Mabuse, the Gambler), where men mesmerized by card playing smoke cigarettes while gambling.
Female smokers in film were also early on associated with a type of sensuous and seductive sexuality, most notably personified by German film star Marlene Dietrich. Similarly, actors like Humphrey Bogart and Audrey Hepburn have been closely identified with their smoker persona, and some of their most famous portraits and roles have involved them being haloed by a mist of cigarette smoke. Hepburn often enhanced the glamor with a cigarette holder, most notably in the film Breakfast at Tiffany's. Smoking could also be used as a means to subvert censorship, as two cigarettes burning unattended in an ashtray were often used to suggest sexual activity.
Since World War II, smoking has gradually become less frequent on screen as the obvious health hazards of smoking have become more widely known. With the anti-smoking movement gaining greater respect and influence, conscious attempts not to show smoking on screen are now undertaken in order to avoid encouraging smoking or giving it positive associations, particularly for family films. Smoking on screen is more common today among characters who are portrayed as anti-social or even criminal.
According to a 2019 study, the introduction of television in the United States led to a substantial increase in smoking, in particular among 16–21-year-olds. The study suggested "that television increased the share of smokers in the population by 5–15 percentage points, generating roughly 11 million additional smokers between 1946 and 1970."
Literature
Just as in other types of fiction, smoking has had an important place in literature and smokers are often portrayed as characters with great individuality, or outright eccentrics, something typically personified in one of the most iconic smoking literary figures of all, Sherlock Holmes. Other than being a frequent part of short stories and novels, smoking has spawned endless eulogies, praising its qualities and affirming the author's identity as a devoted smoker. Especially during the late 19th century and early 20th century, a panoply of books with titles like Tobacco: Its History and associations (1876), Cigarettes in Fact and Fancy (1906) and Pipe and Pouch: The Smokers Own Book of Poetry (1905) were written in the UK and the US. The titles were written by men for other men and contained general tidbits and poetic musings about the love for tobacco and all things related to it, and frequently praised the refined bachelor's life. The Fragrant Weed: Some of the Good Things Which Have been Said or Sung about Tobacco, published in 1907, contained, among many others, the following lines from the poem A Bachelor's Views by Tom Hall that were typical of the attitude of many of the books:
These works were all published in an era before the cigarette had become the dominant form of tobacco consumption and pipes, cigars, and chewing tobacco were still commonplace. Many of the books were published in novel packaging that would attract the learned smoking gentleman. Pipe and Pouch came in a leather bag resembling a tobacco pouch and Cigarettes in Fact and Fancy (1901) came bound in leather, packaged in an imitation cardboard cigar box. By the late 1920s, the publication of this type of literature largely abated and was only sporadically revived in the later 20th century.
Music
There have been few examples of tobacco in music in early modern times, though there are occasional signs of influence in pieces such as Johann Sebastian Bach's Enlightening Thoughts of a Tobacco-Smoker. However, from the early 20th century and onwards smoking has been closely associated with popular music. Jazz was from early on closely intertwined with the smoking that was practiced in the venues where it was played, such as bars, dance halls, jazz clubs and even brothels. The rise of jazz coincided with the expansion of the modern tobacco industry, and in the United States also contributed to the spread of cannabis. The latter went under names like "tea", "muggles" and "reefer" in the jazz community and was so influential in the 1920s and 30s that it found its way into songs composed at the time such as Louis Armstrong's Muggles, Larry Adler's Smoking Reefers, and Don Redman's Chant of The Weed. The popularity of marijuana among jazz musicians remained high until the 1940s and 50s, when it was partially replaced by the use of heroin.
Another form of modern popular music that has been closely associated with cannabis smoking is reggae, a style of music that originated in Jamaica in the late 1950s and early 60s. Cannabis, or ganja, is believed to have been introduced to Jamaica in the mid-19th century by Indian immigrant labor and was primarily associated with Indian workers until it was appropriated by the Rastafari movement in the middle of the 20th century. The Rastafari considered cannabis smoking to be a way to come closer to God, or Jah, an association that was greatly popularized by reggae icons such as Bob Marley and Peter Tosh in the 1960s and 70s.
Economics
Estimates claim that smokers cost the U.S. economy $97.6 billion a year in lost productivity and that an additional $96.7 billion is spent on public and private health care combined. This is over 1% of the gross domestic product. A male smoker in the United States that smokes more than one pack a day can expect an average increase of $19,000 just in medical expenses over the course of his lifetime. A U.S. female smoker that also smokes more than a pack a day can expect an average of $25,800 additional healthcare costs over her lifetime.
| Biology and health sciences | Drugs and medication | null |
13872655 | https://en.wikipedia.org/wiki/Trimerophytopsida | Trimerophytopsida | Trimerophytopsida (or Trimeropsida) is a class of early vascular plants from the Devonian, informally called trimerophytes. It contains genera such as Psilophyton. This group is probably paraphyletic, and is believed to be the ancestral group from which both the ferns and seed plants evolved. Different authors have treated the group at different taxonomic ranks using the names Trimerophyta, Trimerophytophyta, Trimerophytina, Trimerophytophytina and Trimerophytales.
Taxonomy
At first most of the early land plants other than the bryophytes (i.e. the polysporangiophytes) were placed in a single class Psilophyta, established in 1917 by Kidston and Lang. As additional fossils were discovered and described, it became apparent that the Psilophyta were not a homogeneous group of plants. In 1968 Banks first proposed splitting this taxon into three groups, which he put at the rank of subdivision; he clarified his proposal in 1975. One of the three groups was the Trimerophytina. The subdivision is based on the type genus Trimerophyton, which might be expected to produce 'Trimerophytophytina' as the name of the subdivision, but the International Code of Nomenclature for algae, fungi, and plants allows the 'phyton' part of a genus name optionally to be omitted before '-ophyta', '-ophytina' and '-opsida'.
The group has also since been treated as a division under the name Trimerophyta or Trimerophytophyta, as a class under the name Trimeropsida or Trimerophytopsida (as here), and as an order under the name Trimerophytales.
Subphylum †Trimerophytina Banks 1975
Class †Trimerophytopsida Foster & Gifford 1974 [Trimeropsida Banks 1975; Psilophytopsida Kidston & Lang 1917; Psilophytidae Nemejc 1963]
Order †Trimerophytales Banks ex Kasper et al. 1974 [Psilophytales Pia 1924]
Family †Trimerophytaceae Banks 1967 [Psilophytaceae Hirmer 1927; Hostinellaceae Pia 1924; Dawsonitaceae Nemejc 1963]
Genus †Dawsonites Nemejc 1963
Genus †Hostinella Barrande ex Stur 1882
Genus †Oocampsa Andrews, Gensel & Kasper 1975
Genus †Euphyllophyton Hao & Beck
Genus †Pauthecophyton Xue et al. 2012
Genus †Psilophyton Dawson 1859 emend. Hueber & Banks 1967
Genus †Trimerophyton Hopping 1956
| Biology and health sciences | Pteridophytes | Plants |
13873010 | https://en.wikipedia.org/wiki/VizieR | VizieR | The VizieR Catalogue Service is an astronomical catalog service provided by Centre de données astronomiques de Strasbourg (CDS, Strasbourg Data Centre for astronomy). The origin of the service dates back to 1993, when it was founded by the European Space Agency as the European Space Information System (ESIS) Catalogue Browser. Initially intended to serve the space science community, the ESIS project pre-dates the World Wide Web as a network database allowing uniform to a heterogeneous set of catalogues and data.
The CDS had for years collected and disseminated astronomical data, so the original ESIS Catalogue Browser was transferred to and stored there. Since its inception in 1996, VizieR has become a reference point for astronomers worldwide engaged in research, who come to access catalogued data regularly published in astronomical journals. The new VizieR service was refurbished in 1997 by the CDS to better serve the community in terms of searching capabilities and data volume. As of February 2021 it contains more than 20000 catalogues. It is a major data source as part of the Virtual Observatory.
| Physical sciences | Databases | Astronomy |
13873779 | https://en.wikipedia.org/wiki/Sugarcane | Sugarcane | Sugarcane or sugar cane is a species of tall, perennial grass (in the genus Saccharum, tribe Andropogoneae) that is used for sugar production. The plants are 2–6 m (6–20 ft) tall with stout, jointed, fibrous stalks that are rich in sucrose, which accumulates in the stalk internodes. Sugarcanes belong to the grass family, Poaceae, an economically important flowering plant family that includes maize, wheat, rice, and sorghum, and many forage crops. It is native to New Guinea.
Sugarcane was an ancient crop of the Austronesian and Papuan people. The best evidence available today points to the New Guinea area as the site of the original domestication of Saccharum officinarum. It was introduced to Polynesia, Island Melanesia, and Madagascar in prehistoric times via Austronesian sailors. It was also introduced to southern China and India by Austronesian traders around 1200 to 1000 BC. The Persians and Greeks encountered the famous "reeds that produce honey without bees" in India between the sixth and fourth centuries BC. They adopted and then spread sugarcane agriculture. Merchants began to trade in sugar, which was considered a luxurious and expensive spice, from India. In the 18th century, sugarcane plantations began in the Caribbean, South American, Indian Ocean, and Pacific island nations. The need for sugar crop laborers became a major driver of large migrations, some people voluntarily accepting indentured servitude and others forcibly imported as slaves.
Grown in tropical and subtropical regions, sugarcane is the world's largest crop by production quantity, totalling 1.9 billion tonnes in 2020, with Brazil accounting for 40% of the world total. Sugarcane accounts for 79% of sugar produced globally (most of the rest is made from sugar beets). About 70% of the sugar produced comes from Saccharum officinarum and its hybrids. All sugarcane species can interbreed, and the major commercial cultivars are complex hybrids.
White sugar is produced from sugarcane in specialized mill factories. Sugarcane reeds are used to make pens, mats, screens, and thatch. The young, unexpanded flower head of Saccharum edule (duruka) is eaten raw, steamed, or toasted, and prepared in various ways in Southeast Asia, such as certain island communities of Indonesia as well as in Oceanic countries like Fiji. The direct use of sugar cane to produce ethanol for biofuel is projected to potentially surpass the production of white sugar as an end product.
Etymology
The term sugarcane is a combination of two words: "sugar" and "cane". The former ultimately derives from Sanskrit शर्करा (śárkarā) as the crop originated in Southeast Asia. As sugar was traded and spread West, this became سُكَّر (sukkar) in Arabic, saccharum or succarum in Latin, zúcchero in Italian, and eventually sucre in both Middle French and Middle English. The second term "cane" began to be used alongside it as the crop was grown on plantations in the Caribbean. Ganna is the Hindi word for sugarcane.
Characteristics
Sugarcane, a perennial tropical grass, exhibits a unique growth pattern characterized by lateral shoots emerging at its base, leading to the development of multiple stems. These stems typically attain a height of 3 to 4 meters (approximately 10 to 13 feet) and possess a diameter of about 5 centimeters (approximately 2 inches). As these stems mature, they evolve into cane stalks, constituting a substantial portion of the entire plant, accounting for roughly 75% of its composition.
A fully mature cane stalk generally comprises a composition of around 11–16% fiber, 12–16% soluble sugars, 2–3% nonsugar carbohydrates, and 63–73% water content. The successful cultivation of sugarcane hinges on a delicate interplay of several factors, including climatic conditions, soil properties, irrigation methods, fertilization practices, pest and disease management, the selection of specific varieties, and the timing of the harvest.
In terms of yield, the average production of cane stalk stands at 60–70 tonnes per hectare (equivalent to 24–28 long tons per acre or 27–31 short tons per acre) annually. However, this yield figure is not fixed and can vary significantly, ranging from 30 to 180 tonnes per hectare. This variance is contingent upon the level of knowledge applied and the approach to crop management embraced in the cultivation of sugarcane. Ultimately, the successful cultivation of this valuable crop demands a thoughtful integration of various factors to optimize its growth and productivity.
Sugarcane is a cash crop, but it is also used as livestock fodder. Sugarcane genome is one of the most complex plant genomes known, mostly due to interspecific hybridization and polyploidization.
History
The two centers of domestication for sugarcane are one for Saccharum officinarum by Papuans in New Guinea and another for Saccharum sinense by Austronesians in Taiwan and southern China. Papuans and Austronesians originally primarily used sugarcane as food for domesticated pigs. The spread of both S. officinarum and S. sinense is closely linked to the migrations of the Austronesian peoples. Saccharum barberi was only cultivated in India after the introduction of S. officinarum.
S. officinarum was first domesticated in New Guinea and the islands east of the Wallace Line by Papuans, where it is the modern center of diversity. Beginning around 6,000 BP, several strains were selectively bred from the native Saccharum robustum. From New Guinea, it spread westwards to Maritime Southeast Asia after contact with Austronesians, where it hybridized with Saccharum spontaneum.
The second domestication center is southern China and Taiwan, where S. sinense was a primary cultigen of the Austronesian peoples. Words for sugarcane are reconstructed as *təbuS or *CebuS in Proto-Austronesian, which became *tebuh in Proto-Malayo-Polynesian. It was one of the original major crops of the Austronesian peoples from at least 5,500 BP. Introduction of the sweeter S. officinarum may have gradually replaced it throughout its cultivated range in maritime Southeast Asia.
From Insular Southeast Asia, S. officinarum was spread eastward into Polynesia and Micronesia by Austronesian voyagers as a canoe plant by around 3,500 BP. It was also spread westward and northward by around 3,000 BP to China and India by Austronesian traders, where it further hybridized with S. sinense and S. barberi. From there, it spread further into western Eurasia and the Mediterranean.
The earliest known production of crystalline sugar began in northern India. The earliest evidence of sugar production comes from ancient Sanskrit and Pali texts. Around the eighth century, Muslim and Arab traders introduced sugar from medieval India to the other parts of the Abbasid Caliphate in the Mediterranean, Mesopotamia, Egypt, North Africa, and Andalusia. By the 10th century, sources state that every village in Mesopotamia grew sugarcane. It was among the early crops brought to the Americas by the Spanish, mainly Andalusians, from their fields in the Canary Islands, and the Portuguese from their fields in the Madeira Islands. An article on sugarcane cultivation in Spain is included in Ibn al-'Awwam's 12th-century Book on Agriculture.
The first chemically refined sugar appeared on the scene in India about 2,500 years ago. From there, the technique spread east towards China, and west towards Persia and the early Islamic worlds, eventually reaching the Mediterranean in the 13th century. Cyprus and Sicily became important centers for sugar production.
In colonial times, sugar formed one side of the triangle trade of New World raw materials, along with European manufactured goods, and African slaves. Christopher Columbus first brought sugarcane to the Caribbean (and the New World) during his second voyage to the Americas, initially to the island of Hispaniola (modern day Haiti and the Dominican Republic). The first sugar harvest happened in Hispaniola in 1501; many sugar mills were constructed in Cuba and Jamaica by the 1520s. The Portuguese introduced sugarcane to Brazil. By 1540, there were 800 cane sugar mills in Santa Catarina Island and another 2,000 on the north coast of Brazil, Demarara, and Suriname.
Sugar, often in the form of molasses, was shipped from the Caribbean to Europe or New England, where it was used to make rum. The profits from the sale of sugar were then used to purchase manufactured goods, which were then shipped to West Africa, where they were bartered for slaves. The slaves were then brought back to the Caribbean to be sold to sugar planters. The profits from the sale of the slaves were then used to buy more sugar, which was shipped to Europe. Toil in the sugar plantations became a main basis for a vast network of forced population movement, supplying people to work under brutal coercion.
The passage of the 1833 Slavery Abolition Act led to the abolition of slavery through most of the British Empire, and many of the emancipated slaves no longer worked on sugarcane plantations when they had a choice. West Indian planters, therefore, needed new workers, and they found cheap labour in China and India. The people were subject to indenture, a long-established form of contract, which bound them to unfree labour for a fixed term. The conditions where the indentured servants worked were frequently abysmal, owing to a lack of care among the planters. The first ships carrying indentured labourers from India left in 1836. The migrations to serve sugarcane plantations led to a significant number of ethnic Indians, Southeast Asians, and Chinese people settling in various parts of the world. In some islands and countries, the South Asian migrants now constitute between 10 and 50% of the population. Sugarcane plantations and Asian ethnic groups continue to thrive in countries such as Fiji, South Africa, Myanmar, Sri Lanka, Malaysia, Indonesia, the Philippines, Guyana, Jamaica, Trinidad, Martinique, French Guiana, Guadeloupe, Grenada, St. Lucia, St. Vincent, St. Kitts, St. Croix, Suriname, Nevis, and Mauritius.
Between 1863 and 1900, merchants and plantation owners in Queensland and New South Wales (now part of the Commonwealth of Australia) brought between 55,000 and 62,500 people from the South Pacific islands to work on sugarcane plantations. An estimated one-third of these workers were coerced or kidnapped into slavery (known as blackbirding); many others were paid very low wages. Between 1904 and 1908, most of the 10,000 remaining workers were deported in an effort to keep Australia racially homogeneous and protect white workers from cheap foreign labour.
Cuban sugar derived from sugarcane was exported to the USSR, where it received price supports and was ensured a guaranteed market. The 1991 dissolution of the Soviet state forced the closure of most of Cuba's sugar industry.
Sugarcane remains an important part of the economy of Cuba, Guyana, Belize, Barbados, and Haiti, along with the Dominican Republic, Guadeloupe, Jamaica, and other islands.
About 70% of the sugar produced globally comes from S. officinarum and hybrids using this species.
Sugar occupies 26,942,686 hectares of land across the globe and is the third most valuable crop.
Cultivation
Sugarcane cultivation requires a tropical or subtropical climate, with a minimum of of annual moisture. It is one of the most efficient photosynthesizers in the plant kingdom. It is a C4 plant, able to convert up to 1% of incident solar energy into biomass. In primary growing regions across the tropics and subtropics, sugarcane crops can produce over 15 kg/m2 of cane.
Sugar cane accounted for around 21% of the global crop production over the 2000–2021 period. The Americas was the leading region in the production of sugar cane (52% of the world total).
Once a major crop of the southeastern region of the United States, sugarcane cultivation declined there during the late 20th century, and is primarily confined to small plantations in Florida, Louisiana, and southeast Texas in the 21st century. Sugarcane cultivation ceased in Hawaii when the last operating sugar plantation in the state shut down in 2016.
Sugarcane is cultivated in the tropics and subtropics in areas with a plentiful supply of water for a continuous period of more than 6–7 months each year, either from natural rainfall or through irrigation. The crop does not tolerate severe frosts. Therefore, most of the world's sugarcane is grown between 22°N and 22°S, and some up to 33°N and 33°S. When sugarcane crops are found outside this range, such as the Natal region of South Africa, it is normally due to anomalous climatic conditions in the region, such as warm ocean currents that sweep down the coast. In terms of altitude, sugarcane crops are found up to close to the equator in countries such as Colombia, Ecuador, and Peru.
Sugarcane can be grown on many soils ranging from highly fertile, well-drained mollisols, through heavy cracking vertisols, infertile acid oxisols and ultisols, peaty histosols, to rocky andisols. Both plentiful sunshine and water supplies increase cane production. This has made desert countries with good irrigation facilities such as Egypt some of the highest-yielding sugarcane-cultivating regions. Sugarcane consumes 9% of the world's potash fertilizer production.
Although some sugarcanes produce seeds, modern stem cutting has become the most common reproduction method. Each cutting must contain at least one bud, and the cuttings are sometimes hand-planted. In more technologically advanced countries, such as the United States and Australia, billet planting is common. Billets (stalks or stalk sections) harvested by a mechanical harvester are planted by a machine that opens and recloses the ground. Once planted, a stand can be harvested several times; after each harvest, the cane sends up new stalks, called ratoons. Successive harvests give decreasing yields, eventually justifying replanting. Two to 10 harvests are usually made depending on the type of culture. In a country with a mechanical agriculture looking for a high production of large fields, as in North America, sugarcanes are replanted after two or three harvests to avoid a lowering yields. In countries with a more traditional type of agriculture with smaller fields and hand harvesting, as in the French island of Réunion, sugarcane is often harvested up to 10 years before replanting.
Sugarcane is harvested by hand and mechanically. Hand harvesting accounts for more than half of production, and is dominant in the developing world. In hand harvesting, the field is first set on fire. The fire burns up dry leaves, and chases away or kills venomous snakes, without harming the stalks and roots. Harvesters then cut the cane just above ground-level using cane knives or machetes. A skilled harvester can cut of sugarcane per hour.
Mechanical harvesting uses a combine, or sugarcane harvester. The Austoft 7000 series, the original modern harvester design, has now been copied by other companies, including Cameco / John Deere. The machine cuts the cane at the base of the stalk, strips the leaves, chops the cane into consistent lengths and deposits it into a transporter following alongside. The harvester then blows the trash back onto the field. Such machines can harvest each hour, but harvested cane must be rapidly processed. Once cut, sugarcane begins to lose its sugar content, and damage to the cane during mechanical harvesting accelerates this decline. This decline is offset because a modern chopper harvester can complete the harvest faster and more efficiently than hand cutting and loading. Austoft also developed a series of hydraulic high-lift infield transporters to work alongside its harvesters to allow even more rapid transfer of cane to, for example, the nearest railway siding. This mechanical harvesting does not require the field to be set on fire; the residue left in the field by the machine consists of cane tops and dead leaves, which serve as mulch for the next planting.
Pests
The cane beetle (also known as cane grub) can substantially reduce crop yield by eating roots; it can be controlled with imidacloprid (Confidor) or chlorpyrifos (Lorsban). Other important pests are the larvae of some butterfly/moth species, including the turnip moth, the sugarcane borer (Diatraea saccharalis), the African sugarcane borer (Eldana saccharina), the Mexican rice borer (Eoreuma loftini), the African armyworm (Spodoptera exempta), leafcutter ants, termites, spittlebugs (especially Mahanarva fimbriolata and Deois flavopicta), and Migdolus fryanus (a beetle). The planthopper insect Eumetopina flavipes acts as a virus vector, which causes the sugarcane disease ramu stunt. Sesamia grisescens is a major pest in Papua New Guinea and so is a serious concern for the Australian industry were it to cross over. To head off such a problem, the Federal Government has pre-announced that they would cover 80% of response costs if it were necessary.
Pathogens
Numerous pathogens infect sugarcane, such as sugarcane grassy shoot disease caused by Candidatus Phytoplasma sacchari, whiptail disease or sugarcane smut, pokkah boeng caused by Fusarium moniliforme, Xanthomonas axonopodis bacteria causes Gumming Disease, and red rot disease caused by Colletotrichum falcatum. Viral diseases affecting sugarcane include sugarcane mosaic virus, maize streak virus, and sugarcane yellow leaf virus.
Yang et al., 2017 provides a genetic map developed for USDA ARS-run breeding programs for brown rust of sugarcane.
Nitrogen fixation
Some sugarcane varieties are capable of fixing atmospheric nitrogen in association with the bacterium Gluconacetobacter diazotrophicus. Unlike legumes and other nitrogen-fixing plants that form root nodules in the soil in association with bacteria, G. diazotrophicus lives within the intercellular spaces of the sugarcane's stem. Coating seeds with the bacteria was assayed in 2006 with the intention of enabling crop species to fix nitrogen for its own use.
Conditions for sugarcane workers
At least 20,000 people are estimated to have died of chronic kidney disease in Central America in the past two decades, most of them sugarcane workers along the Pacific coast. This may be due to working long hours in the heat without adequate fluid intake.
Additionally, some of the workers are being exposed to hazards such as: high temperatures, harmful pesticides, and poisonous or venomous animals. This occurs during the process of cutting the sugarcane manually, causing physical ailments due to constant repetitive movements for hours every work day.
Processing
Traditionally, sugarcane processing requires two stages. Mills extract raw sugar from freshly harvested cane and "mill-white" sugar is sometimes produced immediately after the first stage at sugar-extraction mills, intended for local consumption. Sugar crystals appear naturally white in color during the crystallization process. Sulfur dioxide is added to inhibit the formation of color-inducing molecules and to stabilize the sugar juices during evaporation. Refineries, often located nearer to consumers in North America, Europe, and Japan, then produce refined white sugar, which is 99% sucrose. These two stages are slowly merging. Increasing affluence in the sugarcane-producing tropics increases demand for refined sugar products, driving a trend toward combined milling and refining.
Milling
Sugarcane processing produces cane sugar (sucrose) from sugarcane. Other products of the processing include bagasse, molasses, and filtercake.
Bagasse, the residual dry fiber of the cane after cane juice has been extracted, is used for several purposes:
fuel for the boilers and kilns
production of paper, paperboard products, and reconstituted panelboard
agricultural mulch
as a raw material for production of chemicals
The primary use of bagasse and bagasse residue is as a fuel source for the boilers in the generation of process steam in sugar plants. Dried filtercake is used as an animal feed supplement, fertilizer, and source of sugarcane wax.
Molasses is produced in two forms: blackstrap, which has a characteristic strong flavor, and a purer molasses syrup. Blackstrap molasses is sold as a food and dietary supplement. It is also a common ingredient in animal feed, and is used to produce ethanol, rum, and citric acid. Purer molasses syrups are sold as molasses, and may also be blended with maple syrup, invert sugars, or corn syrup. Both forms of molasses are used in baking.
Refining
Sugar refining further purifies the raw sugar. It is first mixed with heavy syrup and then centrifuged in a process called "affination". Its purpose is to wash away the sugar crystals' outer coating, which is less pure than the crystal interior. The remaining sugar is then dissolved to make a syrup, about 60% solids by weight.
The sugar solution is clarified by the addition of phosphoric acid and calcium hydroxide, which combine to precipitate calcium phosphate. The calcium phosphate particles entrap some impurities and absorb others, and then float to the top of the tank, where they can be skimmed off. An alternative to this "phosphatation" technique is "carbonatation", which is similar, but uses carbon dioxide and calcium hydroxide to produce a calcium carbonate precipitate.
After filtering any remaining solids, the clarified syrup is decolorized by filtration through activated carbon. Bone char or coal-based activated carbon is traditionally used in this role. Some remaining color-forming impurities are adsorbed by the carbon. The purified syrup is then concentrated to supersaturation and repeatedly crystallized in a vacuum, to produce white refined sugar. As in a sugar mill, the sugar crystals are separated from the molasses by centrifuging. Additional sugar is recovered by blending the remaining syrup with the washings from affination and again crystallizing to produce brown sugar. When no more sugar can be economically recovered, the final molasses still contains 30-35% sucrose and 10–25% glucose and fructose.
To produce granulated sugar, in which individual grains do not clump, sugar must be dried, first by heating in a rotary dryer, and then by blowing cool air through it for several days.
Ribbon cane syrup
Ribbon cane is a subtropical type that was once widely grown in the Southern United States, as far north as coastal North Carolina. The juice was extracted with horse- or mule-powered crushers; the juice was boiled, like maple syrup, in a flat pan, and then used in the syrup form as a food sweetener. It is not currently a commercial crop, but a few growers find ready sales for their product.
Production
In 2022, global production of sugarcane was 1.92 billion tonnes, with Brazil producing 38% of the world total, India with 23%, and China producing 5% (table).
Worldwide, 26 million hectares were devoted to sugarcane cultivation in 2020. The average worldwide yield of sugarcane crops in 2022 was 74 tonnes per hectare, led by Peru with 121 tonnes per hectare. The theoretical possible yield for sugarcane is about 280 tonnes per hectare per year, and small experimental plots in Brazil have demonstrated yields of 236–280 tonnes of cane per hectare.
From 2008 to 2016, production of standards-compliant sugarcane experienced a compound annual growth rate of about 52%, while conventional sugarcane increased at less than 1%.
Environmental impacts
Soil degradation and erosion
The cultivation of sugarcane can lead to increased soil loss through the removal of soil at harvest, as well as improper irrigation practices, which can result in erosion. Erosion is especially significant when the sugarcane is grown on slopes or hillsides, which increases the rate of water runoff. Generally, it is recommended that sugarcane is not planted in areas with a slope greater than 8%. However, in certain areas, such as parts of the Caribbean and South Africa, slopes greater than 20% have been planted. Increased erosion can lead to the removal of organic and nutrient-rich material, which can decrease future crop yields. It can also result in sediments and other pollutants being washed into aquatic habitats, which can result in a wide range of environmental issues, including eutrophication and acidification.
Sugarcane cultivation can also result in soil compaction, which is caused by the use of heavy, infield machinery. Along with impacting invertebrate and fauna within the upper layers of the soil, compaction can also lead to decreased porosity. This in turn can increase surface runoff, resulting in greater leaching and erosion.
Habitat destruction
Due to the large quantity of water required, sugarcane cultivation heavily relies on irrigation. Additionally, since large amounts of soil are removed with the crop during harvest, significant washing occurs during the processing phase. In many countries, such as India and Australia, this requirement has placed a strain on available resources, requiring the construction of barrages and other dams. This has altered the amount of water reaching aquatic habitats, and has contributed to the degradation of ecosystems such as the Great Barrier Reef and Indus Delta.
Sugarcane has also contributed to habitat destruction through the clearance of land. Seven countries around the world devote more than 50% of their land to the cultivation of sugarcane. Sugarcane fields have replaced tropical rain forests and wetlands. While the majority of this clearance occurred in the past, expansions have occurred within the past couple decades, further contributing to habitat destruction.
Mitigation efforts
A wide variety of mitigation efforts can be implemented to reduce the impacts of sugarcane cultivation. Among these efforts is switching to alternative irrigation techniques, such as drip irrigation, which are more water efficient. Water efficiency can also be improved by employing methods such as trash mulching, which has been shown to increase water intake and storage. Along with reducing the overall water use, this method can also decrease soil runoff, and therefore prevent pollutants from entering the environment. In areas with a slope greater than 11%, it is also recommended that zero tillage or cane strip planting are implemented to help prevent soil loss.
Sugarcane processing produces a wide variety of pollutants, including heavy metals and bagasse, which can be released into the environment through wastewater discharge. To prevent this, alternative treatment methods such as high rate anaerobic digestions can be implemented to better treat this wastewater. Stormwater drains can also be installed to prevent uncontrolled runoff from reaching aquatic ecosystems.
Ethanol
Ethanol is generally available as a byproduct of sugar production. It can be used as a biofuel alternative to gasoline, and is widely used in cars in Brazil. It is an alternative to gasoline, and may become the primary product of sugarcane processing, rather than sugar
In Brazil, gasoline is required to contain at least 22% bioethanol. This bioethanol is sourced from Brazil's large sugarcane crop.
The production of ethanol from sugarcane is more energy efficient than from corn or sugar beets or palm/vegetable oils, particularly if cane bagasse is used to produce heat and power for the process. Furthermore, if biofuels are used for crop production and transport, the fossil energy input needed for each ethanol energy unit can be very low. EIA estimates that with an integrated sugarcane to ethanol technology, the well-to-wheels CO2 emissions can be 90% lower than conventional gasoline. A textbook on renewable energy describes the energy transformation:
Presently, 75 tons of raw sugarcane are produced annually per hectare in Brazil. The cane delivered to the processing plant is called burned and cropped (b&c), and represents 77% of the mass of the raw cane. The reason for this reduction is that the stalks are separated from the leaves (which are burned and whose ashes are left in the field as fertilizer), and from the roots that remain in the ground to sprout for the next crop. Average cane production is, therefore, 58 tons of b&c per hectare per year.
Each ton of b&c yields 740 kg of juice (135 kg of sucrose and 605 kg of water) and 260 kg of moist bagasse (130 kg of dry bagasse). Since the lower heating value of sucrose is 16.5 MJ/kg, and that of the bagasse is 19.2 MJ/kg, the total heating value of a ton of b&c is 4.7 GJ of which 2.2 GJ come from the sucrose and 2.5 from the bagasse.
Per hectare per year, the biomass produced corresponds to 0.27 TJ. This is equivalent to 0.86 W per square meter. Assuming an average insolation of 225 W per square meter, the photosynthetic efficiency of sugar cane is 0.38%.
The 135 kg of sucrose found in 1 ton of b&c are transformed into 70 litres of ethanol with a combustion energy of 1.7 GJ. The practical sucrose-ethanol conversion efficiency is, therefore, 76% (compare with the theoretical 97%).
One hectare of sugar cane yields 4,000 litres of ethanol per year (without any additional energy input, because the bagasse produced exceeds the amount needed to distill the final product). This, however, does not include the energy used in tilling, transportation, and so on. Thus, the solar energy-to-ethanol conversion efficiency is 0.13%.
Bagasse applications
Sugarcane is a major crop in many countries. It is one of the plants with the highest bioconversion efficiency. Sugarcane crop is able to efficiently fix solar energy, yielding some 55 tonnes of dry matter per hectare of land annually. After harvest, the crop produces sugar juice and bagasse, the fibrous dry matter. This dry matter is biomass with potential as fuel for energy production. Bagasse can also be used as an alternative source of pulp for paper production.
Sugarcane bagasse is a potentially abundant source of energy for large producers of sugarcane, such as Brazil, India, and China. According to one report, with use of latest technologies, bagasse produced annually in Brazil has the potential of meeting 20% of Brazil's energy consumption by 2020.
Electricity production
A number of countries, in particular those lacking fossil fuels, have implemented energy conservation and efficiency measures to minimize the energy used in cane processing, and export any excess electricity to the grid. Bagasse is usually burned to produce steam, which in turn creates electricity. Current technologies, such as those in use in Mauritius, produce over 100 kWh of electricity per tonne of bagasse. With a total world harvest of over one billion tonnes of sugarcane per year, the global energy potential from bagasse is over 100,000 GWh. Using Mauritius as a reference, an annual potential of 10,000 GWh of additional electricity could be produced throughout Africa. Electrical generation from bagasse could become quite important, particularly to the rural populations of sugarcane producing nations.
Recent cogeneration technology plants are being designed to produce from 200 to over 300 kWh of electricity per tonne of bagasse. As sugarcane is a seasonal crop, shortly after harvest the supply of bagasse would peak, requiring power generation plants to strategically manage the storage of bagasse.
Biogas production
A greener alternative to burning bagasse for the production of electricity is to convert bagasse into biogas. Technologies are being developed to use enzymes to transform bagasse into advanced biofuel and biogas.
Sugarcane as food
In most countries where sugarcane is cultivated, several foods and popular dishes are derived directly from it, such as:
Raw sugarcane: chewed to extract the juice
Sayur nganten: an Indonesian soup made with the stem of trubuk (Saccharum edule), a type of sugarcane
Sugarcane juice: a combination of fresh juice, extracted by hand or small mills, with a touch of lemon and ice to make a popular drink, known variously as air tebu, usacha rass, guarab, guarapa, guarapo, papelón, aseer asab, ganna sharbat, mosto, caldo de cana, or nước mía
Syrup: a traditional sweetener in soft drinks worldwide but now largely supplanted in the US by high fructose corn syrup, which is less expensive because of corn subsidies and sugar tariffs
Molasses: used as a sweetener and a syrup accompanying other foods, such as cheese or cookies
Jaggery: a solidified molasses, known as gur, gud, or gul in modern Indo-Aryan, is traditionally produced by evaporating juice to make a thick sludge, and then cooling and molding it in buckets. Modern production partially freeze dries the juice to reduce caramelization and lighten its color. It is used as sweetener in cooking traditional entrees, sweets, and desserts.
Falernum: a sweet, and slightly alcoholic drink made from sugarcane juice
Cachaça: the most popular distilled alcoholic beverage in Brazil; it is a liquor made of the distillation of sugarcane juice.
Rum is a liquor made from sugarcane products, typically molasses, but sometimes also cane juice. It is most commonly produced in the Caribbean and environs.
Basi is a fermented alcoholic beverage made from sugarcane juice produced in the Philippines and Guyana.
Panela, solid pieces of sucrose and fructose obtained from the boiling and evaporation of sugarcane juice, is a food staple in Colombia and other countries in South and Central America.
Rapadura is a sweet flour that is one of the simplest refinings of sugarcane juice, common in Latin American countries such as Brazil, Argentina, and Venezuela (where it is known as papelón) and the Caribbean.
Rock candy: crystallized cane juice
Gâteau de Sirop
Viche, a homebrewed Colombian alcoholic beverage
In the 21st century, it is estimated that sugar is potentially responsible for approximately 20% of the caloric content of modern diets.
Sugarcane as feed
Many parts of the sugarcane are commonly used as animal feeds where the plants are cultivated. The leaves make a good forage for ruminants.
Gallery
| Biology and health sciences | Poales | null |
13878118 | https://en.wikipedia.org/wiki/Iron%28III%29%20sulfate | Iron(III) sulfate | Iron(III) sulfate (or ferric sulfate), is a family of inorganic compounds with the formula Fe2(SO4)3(H2O)n. A variety of hydrates are known, including the most commonly encountered form of "ferric sulfate". Solutions are used in dyeing as a mordant, and as a coagulant for industrial wastes. Solutions of ferric sulfate are also used in the processing of aluminum and steel.
Speciation
The various crystalline forms of Fe2(SO4)3(H2O)n are well-defined, often by X-ray crystallography. The nature of the aqueous solutions is often less certain, but aquo-hydroxo complexes such as [Fe(H2O)6]3+ and [Fe(H2O)5(OH)]2+ are often assumed. Regardless, all such solids and solutions feature ferric ions, each with five unpaired electrons. By virtue of this high spin d5 electronic configuration, these ions are paramagnetic and are weak chromophores.
Production
Ferric sulfate solutions are usually generated from iron wastes. The actual identity of the iron species is often vague, but many applications do not demand high purity materials. It is produced on a large scale by treating sulfuric acid, a hot solution of ferrous sulfate, and an oxidizing agent. Typical oxidizing agents include chlorine, nitric acid, and hydrogen peroxide.
Natural occurrences
Iron sulfates occur as a variety of rare (commercially unimportant) minerals. Mikasaite, a mixed iron-aluminium sulfate of chemical formula (Fe3+, Al3+)2(SO4)3 is the name of mineralogical form of iron(III) sulfate. This anhydrous form occurs very rarely and is connected with coal fires. The hydrates are more common, with coquimbite (nonahydrate) as probably the most often met among them. Paracoquimbite is the other, rarely encountered natural nonahydrate. Kornelite (heptahydrate) and quenstedtite (decahydrate) are rarely found. Andradite garnet is a yellow-green example found in Italy. Lausenite (hexa- or pentahydrate) is a doubtful species. All the mentioned natural hydrates are unstable connected with the weathering (aerobic oxidation) of Fe-bearing primary minerals (mainly pyrite and marcasite).
| Physical sciences | Sulfuric oxyanions | Chemistry |
3231565 | https://en.wikipedia.org/wiki/Potassium%20chromate | Potassium chromate | Potassium chromate is the inorganic compound with the formula K2CrO4. This yellow solid is the potassium salt of the chromate anion. It is a common laboratory chemical, whereas sodium chromate is important industrially.
Structure
Two crystalline forms are known, both being very similar to the corresponding potassium sulfate. Orthorhombic β-K2CrO4 is the common form, but it converts to an α-form above 666 °C. These structures are complex, although the chromate ion adopts the typical tetrahedral geometry.
Production and reactions
It is prepared by treating potassium dichromate with potassium hydroxide:
K2Cr2O7(aq) + 2KOH -> 2K2CrO4 + H2O
Or, the fusion of potassium hydroxide and chromium trioxide:
2KOH + CrO3 -> K2CrO4 + H2O
In solution, the behavior of potassium and sodium dichromates are very similar. When treated with lead(II) nitrate, it gives an orange-yellow precipitate, lead(II) chromate.
Applications
Unlike the less expensive sodium salt, potassium salt is mainly used for laboratory work in situations where an anhydrous salt is required. It is as an oxidizing agent in organic synthesis. It is used in qualitative inorganic analysis, e.g. as a colorimetric test for silver ion. It is also used as an indicator in precipitation titrations with silver nitrate and sodium chloride (they can be used as standard as well as titrant for each other) as potassium chromate turns red in the presence of excess of silver ions.
Occurrence
Tarapacaite is the natural, mineral form of potassium chromate. It occurs very rarely and until now is known from only few localities on Atacama Desert.
Safety
As with other Cr(VI) compounds, potassium chromate is carcinogenic. The compound is also corrosive and exposure may produce severe eye damage or blindness. Human exposure further encompasses impaired fertility, heritable genetic damage and harm to unborn children.
| Physical sciences | Metallic oxyanions | Chemistry |
3232420 | https://en.wikipedia.org/wiki/Hydraulic%20cylinder | Hydraulic cylinder | A hydraulic cylinder (also called a linear hydraulic motor) is a mechanical actuator that is used to give a unidirectional force through a unidirectional stroke. It has many applications, notably in construction equipment (engineering vehicles), manufacturing machinery, elevators, and civil engineering.
A hydraulic cylinder is a hydraulic actuator that provides linear motion when hydraulic energy is converted into mechanical movement. It can be likened to a muscle in that, when the hydraulic system of a machine is activated, the cylinder is responsible for providing the motion.
Operation
Hydraulic cylinders get their power from pressurized hydraulic fluid, which is incompressible. Typically oil is used as hydraulic fluid. The hydraulic cylinder consists of a cylinder barrel, in which a piston connected to a piston rod moves back and forth. The barrel is closed on one end by the cylinder bottom (also called the cap) and the other end by the cylinder head (also called the gland) where the piston rod comes out of the cylinder. The piston has sliding rings and seals. The piston divides the inside of the cylinder into two chambers, the bottom chamber (cap end) and the piston rod side chamber (rod end/head-end).
Flanges, trunnions, clevises, and lugs are common cylinder mounting options. The piston rod also has mounting attachments to connect the cylinder to the object or machine component that it is pushing or pulling.
A hydraulic cylinder is the actuator or "motor" side of this system. The "generator" side of the hydraulic system is the hydraulic pump which delivers a fixed or regulated flow of oil to the hydraulic cylinder, to move the piston. There are three types of pump widely used: hydraulic hand pump, hydraulic air pump, and hydraulic electric pump. The piston pushes the oil in the other chamber back to the reservoir. If we assume that the oil enters from the cap end, during extension stroke, and the oil pressure in the rod end/head end is approximately zero, the force F on the piston rod equals the pressure P in the cylinder times the piston area A:
Retraction force difference
For double-acting single-rod cylinders, when the input and output pressures are reversed, there is a force difference between the two sides of the piston due to one side of the piston being covered by the rod attached to it. The cylinder rod reduces the surface area of the piston and reduces the force that can be applied for the retraction stroke.
During the retraction stroke, if the oil is pumped into the head (or gland) at the rod end and the oil from the cap end flows back to the reservoir without pressure, the fluid pressure in the rod end is (Pull Force) / (piston area - piston rod area):
where P is the fluid pressure, Fp is the pulling force, Ap is the piston face area and Ar is the rod cross-section area.
For double-acting, double-rod cylinders, when the piston surface area is equally covered by a rod of equal size on both sides of the head, there is no force difference. Such cylinders typically have their cylinder body affixed to a stationary mount.
Applications
Hydraulic cylinders can be used in any machine where high forces are required, one of the most familiar being earth-moving equipment such as excavators, back hoes and tractors to lift or lower the boom, arm, or bucket. Manufacturing is another popular application where they can be found in hydraulic bending machines, metal sheet shearing machines, particle board or plywood making hot press.
Parts
A hydraulic cylinder has the following parts:
Cylinder barrel
The main function of the cylinder body is to contain cylinder pressure. The cylinder barrel is mostly made from honed tubes. Honed tubes are produced from Suitable To Hone Steel Cold Drawn Seamless Tubes (CDS tubes) or Drawn Over Mandrel (DOM) tubes. Honed tubing is ready to use for hydraulic cylinders without further ID processing. The surface finish of the cylinder barrel is typically 4 to 16 microinch. Honing process and Skiving & Roller burnishing (SRB) process are the two main types of processes for manufacturing cylinder tubes. The piston reciprocates in the cylinder. The cylinder barrel has features of smooth inside surface, high precision tolerance, durable in use, etc.
Cylinder base or cap
The main function of the cap is to enclose the pressure chamber at one end. The cap is connected to the body by means of welding, threading, bolts, or tie rods. Caps also perform as cylinder mounting components [cap flange, cap trunnion, cap clevis]. Capsize is determined based on the bending stress. A static seal / o-ring is used in between cap and barrel (except welded construction).
Cylinder head
The main function of the head is to enclose the pressure chamber from the other end. The head contains an integrated rod sealing arrangement or the option to accept a seal gland. The head is connected to the body by means of threading, bolts, or tie rods. A static seal / o-ring is used in between head and barrel.
Piston
The main function of the piston is to separate the pressure zones inside the barrel. The piston is machined with grooves to fit elastomeric or metal seals and bearing elements. These seals can be single-acting or double-acting. The difference in pressure between the two sides of the piston causes the cylinder to extend and retract. The piston is attached to the piston rod by means of threads, bolts, or nuts to transfer the linear motion.
Piston rod
The piston rod is typically a hard chrome-plated piece of cold-rolled steel that attaches to the piston and extends from the cylinder through the rod-end head. In double rod-end cylinders, the actuator has a rod extending from both sides of the piston and out both ends of the barrel. The piston rod connects the hydraulic actuator to the machine component doing the work. This connection can be in the form of a machine thread or a mounting attachment. The piston rod is highly ground and polished so as to provide a reliable seal and prevent leakage.
Seal gland
The cylinder head is fitted with seals to prevent the pressurized oil from leaking past the interface between the rod and the head. This area is called the seal gland. The advantage of a seal gland is easy removal and seal replacement. The seal gland contains a primary seal, a secondary seal/buffer seal, bearing elements, a wiper/scraper, and a static seal. In some cases, especially in small hydraulic cylinders, the rod gland and the bearing elements are made from a single integral machined part.
Seals
The seals are considered/designed to withstand maximum cylinder working pressure, cylinder speed, operating temperature, working medium, and application. Piston seals are dynamic seals, and they can be single-acting or double-acting. Generally speaking, Elastomer seals made from nitrile rubber, Polyurethane, or other materials are best in lower temperature environments, while seals made of Fluorocarbon Viton are better for higher temperatures. Metallic seals are also available and commonly used cast iron for the seal material. Rod seals are dynamic seals and generally are single-acting. The compounds of rod seals are nitrile rubber, Polyurethane, or Fluorocarbon Viton. Wipers/scrapers are used to eliminate contaminants such as moisture, dirt, and dust, which can cause extensive damage to cylinder walls, rods, seals, and other components. The common compound for wipers is polyurethane. Metallic scrapers are used for sub-zero temperature applications and applications where foreign materials can deposit on the rod. The bearing elements/wear bands are used to eliminate metal to metal contact. The wear bands are designed to withstand maximum side loads. The primary compounds used for wear bands are filled PTFE, woven fabric reinforced polyester resin, and bronze
Other parts
There are many component parts that make up the internal portion of a hydraulic cylinder. All of these pieces combine to create a fully functioning component.
Cylinder base connection
Cushions
Internal Threaded Ductile Heads
Head Glands
Polypak Pistons
Cylinder Head Caps
Butt Plates
Eye Brackets/Clevis Brackets
MP Detachable Mounts
Rod Eyes/Rod Clevis
Pivot Pins
Spherical Ball Bushings
Spherical Rod Eye
Alignment Coupler
Ports and Fittings
Single acting vs. double acting
Single-acting cylinders are economical and the simplest design. Hydraulic fluid enters through a port at one end of the cylinder, which extends the rod by means of area difference. An external force, internal retraction spring or gravity returns the piston rod.
Double acting cylinders have a port at each end or side of the piston, supplied with hydraulic fluid for both the retraction and extension.
Designs
There are primarily two main styles of hydraulic cylinder construction used in the industry: tie rod-style cylinders and welded body-style cylinders.
Tie rod cylinder
Tie rod style hydraulic cylinders use high strength threaded steel rods to hold the two end caps to the cylinder barrel. They are most often seen in industrial factory applications. Small-bore cylinders usually have 4 tie rods, and large bore cylinders may require as many as 16 or 20 tie rods in order to retain the end caps under the tremendous forces produced. Tie rod style cylinders can be completely disassembled for service and repair, and they are not always customizable.
The National Fluid Power Association (NFPA) has standardized the dimensions of hydraulic tie-rod cylinders. This enables cylinders from different manufacturers to interchange within the same mountings.
Welded body cylinder
Welded body cylinders have no tie rods. The barrel is welded directly to the end caps. The ports are welded to the barrel. The front rod gland is usually threaded into or bolted to the cylinder barrel. That allows the piston rod assembly and the rod seals to be removed for service.
Welded body cylinders have a number of advantages over tie rod-style cylinders. Welded cylinders have a narrower body and often a shorter overall length enabling them to fit better into the tight confines of machinery. Welded cylinders do not suffer from failure due to tie rod stretch at high pressures and long strokes. The welded design also lends itself to customization. Special features are easily added to the cylinder body, including special ports, custom mounts, valve manifolds, and so on.
The smooth outer body of welded cylinders also enables the design of multi-stage telescopic cylinders.
Welded body hydraulic cylinders dominate the mobile hydraulic equipment market such as construction equipment (excavators, bulldozers, and road graders) and material handling equipment (forklift trucks, telehandlers, and lift-gates). They are also used by heavy industry in cranes, oil rigs, and large off-road vehicles for above-ground mining operations.
Piston rod construction
The piston rod of a hydraulic cylinder operates both inside and outside the barrel, and consequently both in and out of the hydraulic fluid and surrounding atmosphere.
Coatings
Wear and corrosion-resistant surfaces are desirable on the outer diameter of the piston rod. The surfaces are often applied using coating techniques such as Chrome (Nickel) Plating, Lunac 2+ duplex, Laser Cladding, PTA welding and Thermal Spraying. These coatings can be finished to the desirable surface roughness (Ra, Rz) where the seals give optimum performance. All these coating methods have their specific advantages and disadvantages. It is for this reason that coating experts play a crucial role in selecting the optimum surface treatment procedure for protecting Hydraulic Cylinders.
Cylinders are used in different operational conditions and that makes it a challenge to find the right coating solution. In dredging there might be impact from stones or other parts, in saltwater environments, there are extreme corrosion attacks, in off-shore cylinders facing bending and impact in combination with salt water, and in the steel industry, there are high temperatures involved, etc. There is no single coating solution that successfully combats all the specific operational wear conditions. Every technique has its own benefits and disadvantages.
Length
Piston rods are generally available in lengths that are cut to suit the application. As the common rods have a soft or mild steel core, their ends can be welded or machined for a screw thread.
Distribution of forces on components
The forces on the piston face and the piston head retainer vary depending on which piston head retention system is used.
If a circlip (or any non-preloaded system) is used, the force acting to separate the piston head and the cylinder shaft shoulder is the applied pressure multiplied by the area of the piston head. The piston head and shaft shoulder will separate and the load is fully reacted by the piston head retainer.
If a preloaded system is used the force between the cylinder shaft and piston head is initially the piston head retainer preload value. Once pressure has applied this force will reduce. The piston head and cylinder shaft shoulder will remain in contact unless the applied pressure multiplied by the piston head area exceeds the preload.
The maximum force the piston head retainer will see is the larger of the preload and the applied pressure multiplied by the full piston head area. The load on the piston head retainer is greater than the external load, which is due to the reduced shaft size passing through the piston head. Increasing this portion of shaft reduces the load on the retainer.
Side loading
Side loading is unequal pressure that is not centered on the cylinder rod. This off-center strain can lead to bending of the rod in extreme cases, but more commonly causes leaking due to warping the circular seals into an oval shape. It can also damage and enlarge the bore hole around the rod and the inner cylinder wall around the piston head, if the rod is pressed hard enough sideways to fully compress and deform the seals to make metal-on-metal scraping contact.
The strain of side loading can be directly reduced with the use of internal stop tubes which reduce the maximum extension length, leaving some distance between the piston and bore seal, and increasing leverage to resist warping of the seals. Double pistons also spread out the forces of side loading while also reducing stroke length. Alternately, external sliding guides and hinges can support the load and reduce side loading forces applied directly on the cylinder.
Cylinder mounting methods
Mounting methods also play an important role in cylinder performance. Generally, fixed mounts on the centerline of the cylinder are best for straight-line force transfer and avoiding wear. Common types of mounting include:
Flange mounts—Very strong and rigid, but have little tolerance for misalignment. Experts recommend cap end mounts for thrust loads and rod end mounts where major loading puts the piston rod in tension. Three types are head rectangular flange, head square flange or rectangular head. Flange mounts function optimally when the mounting face attaches to a machine support member.
Side-mounted cylinders—Easy to install and service, but the mounts produce a turning moment as the cylinder applies force to a load, increasing wear and tear. To avoid this, specify a stroke at least as long as the bore size for side mount cylinders (heavy loading tends to make short stroke, large bore cylinders unstable). Side mounts need to be well aligned and the load supported and guided.
Centerline lug mounts —Absorb forces on the centerline, and require dowel pins to secure the lugs to prevent movement at higher pressures or under shock conditions. Dowel pins hold it to the machine when operating at high pressure or under shock loading.
Pivot mounts —Absorb force on the cylinder centerline and let the cylinder change alignment in one plane. Common types include clevises, trunnion mounts and spherical bearings. Because these mounts allow a cylinder to pivot, they should be used with rod-end attachments that also pivot. Clevis mounts can be used in any orientation and are generally recommended for short strokes and small- to medium-bore cylinders.
Special hydraulic cylinders
Telescopic cylinder
The length of a hydraulic cylinder is the total of the stroke, the thickness of the piston, the thickness of bottom and head and the length of the connections. Often this length does not fit in the machine. In that case the piston rod is also used as a piston barrel and a second piston rod is used. These kinds of cylinders are called telescopic cylinders. If we call a normal rod cylinder single stage, telescopic cylinders are multi-stage units of two, three, four, five, or more stages. In general telescopic cylinders are much more expensive than normal cylinders. Most telescopic cylinders are single acting (push). Double acting telescopic cylinders must be specially designed and manufactured.
Plunger cylinder
A hydraulic cylinder without a piston or with a piston without seals is called a plunger cylinder. A plunger cylinder can only be used as a pushing cylinder; the maximum force is piston rod area multiplied by pressure. This means that a plunger cylinder in general has a relatively thick piston rod.
Differential cylinder
A differential cylinder acts like a normal cylinder when pulling. If the cylinder however has to push, the oil from the piston rod side of the cylinder is not returned to the reservoir but goes to the bottom side of the cylinder. In such a way, the cylinder goes much faster, but the maximum force the cylinder can give is like a plunger cylinder. A differential cylinder can be manufactured like a normal cylinder, and only a special control is added.
The above differential cylinder is also called a regenerative cylinder control circuit. This term means that the cylinder is a single rod, double-acting hydraulic cylinder. The control circuit includes a valve and piping which during the extension of the piston, conducts the oil from the rod side of the piston to the other side of the piston instead of to the pump’s reservoir. The oil which is conducted to the other side of the piston is referred to as the regenerative oil.
Position sensing "smart" hydraulic cylinder
Position sensing hydraulic cylinders eliminate the need for a hollow cylinder rod. Instead, an external sensing "bar" using Hall Effect technology senses the position of the cylinder’s piston. This is accomplished by the placement of a permanent magnet within the piston. The magnet propagates a magnetic field through the steel wall of the cylinder, providing a locating signal to the sensor.
| Technology | Hydraulics and pneumatics | null |
3233443 | https://en.wikipedia.org/wiki/Agalychnis | Agalychnis | Agalychnis is a genus of tree frogs native to forests in Mexico, Central America and northwestern South America.
Taxonomy
The following species are recognised in the genus Agalychnis:
External links
| Biology and health sciences | Frogs and toads | Animals |
1683536 | https://en.wikipedia.org/wiki/Superbubble | Superbubble | In astronomy a superbubble or supershell is a cavity which is hundreds of light years across and is populated with hot (106 K) gas atoms, less dense than the surrounding interstellar medium, blown against that medium and carved out by multiple supernovae and stellar winds. The winds, passage and gravity of newly born stars strip superbubbles of any other dust or gas. The Solar System lies near the center of an old superbubble, known as the Local Bubble, whose boundaries can be traced by a sudden rise in dust extinction of exterior stars at distances greater than a few hundred light years.
Formation
The most massive stars, with masses ranging from eight to roughly one hundred solar masses and spectral types of O and early B, are usually found in groups called OB associations. Massive O stars have strong stellar winds, and most of these stars explode as supernovae at the end of their lives.
The strongest stellar winds release kinetic energy of 1051 ergs (1044 J) over the lifetime of a star, which is equivalent to a supernova explosion. These winds can form stellar wind bubbles dozens of light years across.
Inside OB associations, the stars are close enough that their wind bubbles merge, forming a giant bubble called a superbubble.
When stars die, supernova explosions, similarly, drive blast waves that can reach even larger sizes, with expansion velocities up to several hundred km s−1. Stars in OB associations are not gravitationally bound, but they drift apart at small speeds (of around 20 km s−1), and they exhaust their fuel rapidly (after a few millions of years). As a result, most of their supernova explosions occur within the cavity formed by the stellar wind bubbles. These explosions never form a visible supernova remnant, but instead expend their energy in the hot interior as sound waves. Both stellar winds and stellar explosions thus power the expansion of the superbubble in the interstellar medium.
The interstellar gas swept up by superbubbles generally cools, forming a dense shell around the cavity. These shells were first observed in line emission at twenty-one centimeters from hydrogen, leading to the formulation of the theory of superbubble formation. They are also observed in X-ray emission from their hot interiors, in optical line emission from their ionized shells, and in infrared continuum emission from dust swept up in their shells. X-ray and visible emission are typically observed from younger superbubbles, while older, larger objects seen in twenty-one centimeters may even result from multiple superbubbles combining, and so are sometimes distinguished by calling them supershells.
Large enough superbubbles can blow through the entire galactic disk, releasing their energy into the surrounding galactic halo or even into the intergalactic medium.
Examples
LHA 120-N 44 (N44) in the Large Magellanic Cloud.
Anticenter shell, a supershell once called "Snickers"
Henize 70
Monogem Ring
Ophiuchus Superbubble
The Scutum Supershell
Orion-Eridanus Superbubble
The Perseus-Taurus Shell
The Local Bubble
Image gallery
| Physical sciences | Basics_2 | Astronomy |
2345679 | https://en.wikipedia.org/wiki/Priacanthidae | Priacanthidae | The Priacanthidae, the bigeyes, are a family of 18 species of marine ray-finned fishes. "Catalufa" is an alternate common name for some members of the Priacanthidae. The etymology of the scientific name (, to bite + , thorn) refers to the family's very rough, spined scales. The common name of "bigeye" refers to the member species' unusually large eyes, suited to their carnivorous and nocturnal lifestyles. Priacanthidae are typically colored bright red, but some have patterns in silver, dusky brown, or black. Most species reach a maximum total length of about , although in a few species lengths of over are known.
Most members of this family are native to tropical and subtropical parts of the Indian and Pacific Oceans, but four species (Cookeolus japonicus, Heteropriacanthus cruentatus, Priacanthus arenatus, and Pristigenys alta) are found in the Atlantic. They tend to live near rock outcroppings or reefs, although a few are known to inhabit open waters. Many species are found in relatively deep waters, below depths reachable by normal scuba diving. Some species are fished for food.
The earliest identified Priacanthidae fossils date to the middle Eocene epoch of the lower Tertiary period, or roughly 40 to 50 million years ago.
Species
The 18 species in four genera are:
Genus Cookeolus Fowler, 1928
Cookeolus japonicus (Cuvier, 1829) - Longfinned bullseye
†Cookeolus spinolacrymatus Kon & Yoshino, 1997
Genus Heteropriacanthus Fitch & Crooke, 1984
Heteropriacanthus cruentatus (Lacépède, 1801) - Glasseye
Genus Priacanthus Oken, 1817
Priacanthus alalaua Jordan & Evermann, 1903 - Alalaua
Priacanthus arenatus Cuvier, 1829 - Atlantic bigeye
Priacanthus blochii Bleeker, 1853 - Paeony bulleye
Priacanthus fitchi Starnes, 1988
Priacanthus hamrur (Forsskål, 1775) - Moontail bullseye
Priacanthus macracanthus Cuvier, 1829 - Red bigeye
Priacanthus meeki Jenkins, 1903 - Hawaiian bigeye
Priacanthus nasca Starnes, 1988
Priacanthus prolixus Starnes, 1988 - Elongate bulleye
Priacanthus sagittarius Starnes, 1988 - Arrow bulleye
Priacanthus tayenus Richardson, 1846 - Purple-spotted bigeye
Priacanthus zaiserae Starnes & Moyer, 1988
† Priacanthus liui Tao, 1993
Genus Pristigenys Agassiz 1835
Pristigenys alta (Gill, 1862) - Short bigeye
Pristigenys meyeri (Günther, 1872)
Pristigenys niphonia (Cuvier, 1829) - Japanese bigeye
Pristigenys serrula (Gilbert, 1891) - Popeye catalufa
Pristigenys substriatus (Blainville, 1818)
Timeline of genera
| Biology and health sciences | Acanthomorpha | Animals |
2346926 | https://en.wikipedia.org/wiki/Eohippus | Eohippus | Eohippus is an extinct genus of small equid ungulates. The only species is E. angustidens, which was long considered a species of Hyracotherium (now strictly defined as a member of the Palaeotheriidae rather than the Equidae). Its remains have been identified in North America and date to the Early Eocene (Ypresian stage).
Discovery
In 1876, Othniel C. Marsh described a skeleton as Eohippus validus, from (, 'dawn') and (, 'horse'), meaning 'dawn horse'. Its similarities with fossils described by Richard Owen were formally pointed out in a 1932 paper by Clive Forster Cooper. E. validus was moved to the genus Hyracotherium, which had priority as the name for the genus, with Eohippus becoming a junior synonym of that genus. Hyracotherium was recently found to be a paraphyletic group of species, and the genus now includes only H. leporinum. E. validus was found to be identical to an earlier-named species, Orohippus angustidens Cope, 1875, and the resulting binomial is thus Eohippus angustidens.
Description
Eohippus stood at about , or three hands tall, at the shoulder. It has four toes on its front feet and three toes on the hind feet, each toe ending in a hoof. Its incisors, molars and premolars resemble modern Equus. However, a differentiating trait of Eohippus is its large canine teeth.
Stephen Jay Gould comments
In his 1991 essay, "The Case of the Creeping Fox Terrier Clone", Stephen Jay Gould lamented the prevalence of a much-repeated phrase to indicate Eohippus size ("the size of a small Fox Terrier"), even though most readers would be quite unfamiliar with that breed of dog. He concluded that the phrase had its origin in a widely distributed pamphlet by Henry Fairfield Osborn, and proposed that Osborn, a keen fox hunter, could have made a natural association between his horses and the dogs that accompanied them.
| Biology and health sciences | Equidae | Animals |
2347254 | https://en.wikipedia.org/wiki/Cygnus%20A | Cygnus A | Cygnus A (3C 405) is a radio galaxy, one of the strongest radio sources in the sky. A concentrated radio source in Cygnus was discovered by Grote Reber in 1939. In 1946 Stanley Hey and his colleague James Phillips identified that the source scintillated rapidly, and must therefore be a compact object. In 1951, Cygnus A, along with Cassiopeia A, and Puppis A were the first "radio stars" identified with an optical source. Of these, Cygnus A became the first radio galaxy, the other two being nebulae inside the Milky Way. In 1953 Roger Jennison and M K Das Gupta showed it to be a double source. Like all radio galaxies, it contains an active galactic nucleus. The supermassive black hole at the core has a mass of .
Images of the galaxy in the radio portion of the electromagnetic spectrum show two jets protruding in opposite directions from the galaxy's center. These jets extend many times the width of the portion of the host galaxy which emits radiation at visible wavelengths. At the ends of the jets are two lobes with "hot spots" of more intense radiation at their edges. These hot spots are formed when material from the jets collides with the surrounding intergalactic medium.
In 2016, a radio transient was discovered 460 parsecs away from the center of Cygnus A. Between 1989 and 2016, the object, cospatial with a previously-known infrared source, exhibited at least an eightfold increase in radio flux density, with comparable luminosity to the brightest known supernova. Due to the lack of measurements in the intervening years, the rate of brightening is unknown, but the object has remained at a relatively constant flux density since its discovery. The data are consistent with a second supermassive black hole orbiting the primary object, with the secondary having undergone a rapid accretion rate increase. The inferred orbital timescale is of the same order as the activity of the primary source, suggesting the secondary may be perturbing the primary and causing the outflows.
| Physical sciences | Notable galaxies | Astronomy |
2348284 | https://en.wikipedia.org/wiki/Sperrylite | Sperrylite | Sperrylite is a platinum arsenide mineral with the chemical formula and is an opaque metallic tin white mineral which crystallizes in the isometric system with the pyrite group structure. It forms cubic, octahedral or pyritohedral crystals in addition to massive and reniform habits. It has a Mohs hardness of 6–7 and a very high specific gravity of 10.6.
It was discovered by the American chemist Francis Louis Sperry in 1888-89 in ore from the Vermillion Mine, part of the Sudbury Basin discoveries, as he did a fire assay to determine the gold content of the ore for the Canadian Copper Company who were interested by the property.
The most important occurrence of sperrylite is in the nickel ore deposit of Sudbury Basin in Ontario, Canada. It also occurs in the layered igneous complex of the Bushveld region of South Africa and the Oktyabr'skoye copper-nickel deposit near Noril'sk, Russia.
Geologic occurrence
Sperrylite is the most common platinum mineral, it generally occurs with a wide array of other unusual minerals, including cooperite [], laurite [], kotulskite [], merenskyite [], iridium-osmium (Ir-Os) alloys, sudburyite [], omeiite [], testibiopalladite [], and niggliite [], to name a few. It does not readily decompose through normal weathering processes and, consequently, has been reported in widely scattered alluvial deposits.
It was first found as tiny crystals found with rhodolite garnet and corundum during alluvial gem mining in streams draining Mason Mountain, Macon County, North Carolina (Hidden 1898). Sperrylite has been identified in Finland from sulfide deposits generally associated with layered mafic-ultramafic complexes.
Structure
Sperrylite belongs to the pyrite group of minerals and therefore it shares similar structure and crystal habits with them. Analyses typically show minor amounts of rhodium. Trace copper, iron, and antimony as well as intergrowths with Pt-Fe are reported from some occurrences. Sperrylite crystallizes in Pa3, with a =5.9681(l) A. (Szymański, 1979). It has very similar crystal structure as in platarsite (ideally PtAsS). Sperrylite crystals vary considerably in shape and size and are usually enclosed in a variety of host minerals. They are usually closely associated with basemetal sulfide. They are commonly at the edge and partially enclosed by pentlandite, pyrrhotite or chalcopyrite. Seabrook (2004).
Sperrylite is composed of loose aggregate of bright silver cubes, some with octahedral modifications. The grains are mostly anhedral, but a few euhedral grains could also be encountered. Sperrylite is formed by contact metamorphism, as in indicated by the development of triple point annealing contacts with pyrrhotite grains. The grains of sperrylite are surrounded by later veins of pyrite.
Sperrylite is cubic (2/m3) and is typically seen in well-developed cubes or cuboctahedra, some of which are so highly modified that crystal edges and comers appear rounded. (Nicol and Goldschmidt 1903) identified seventeen crystal forms exhibited by sperrylite, including four different trapezohedra, a trisoctahedron, five pyritohedra, and four diploids. Crystals to 2.5 cm have been reported.
Physical properties
Sperrylite is a tin-white mineral known for its brilliant metallic luster, with a grey to black streak. It has indistinct cleavage on {001} and a conchoidal fracture and is brittle. Its hardness is between 6 and 7, and it is quite dense with a calculated specific gravity of 10.78. It has an isometric crystal structure, conchoidal fracture, is non-magnetic and non-radioactive.
Discovery
Sperrylite was first described by H. H. Wells in 1889 from material collected at the Vermilion mine in what is now the Sudbury district, Ontario, Canada. He named it for Francis L. Sperry, chief chemist with the Canadian Copper Company of Sudbury, who collected the original material in 1887 (Mitchell 1985). It occurred in weathered material with colorless transparent cassiterite [], which is thought to have been derived from the oxidation of stannite [].
| Physical sciences | Minerals | Earth science |
2348548 | https://en.wikipedia.org/wiki/Cardboard%20box | Cardboard box | Cardboard boxes are industrially prefabricated boxes, primarily used for packaging goods and materials. Specialists in industry seldom use the term cardboard because it does not denote a specific material. The term cardboard may refer to a variety of heavy paper-like materials,
including card stock, corrugated fiberboard,
and paperboard. Cardboard boxes can be readily recycled.
Terminology
Several types of containers are sometimes called cardboard boxes:
In business and industry, material producers, container manufacturers, packaging engineers,
and standards organizations,
try to use more specific terminology. There is still not complete and uniform usage. Often the term "cardboard" is avoided because it does not define any particular material.
Broad divisions of paper-based packaging materials are:
Paper is thin material mainly used for writing upon, printing upon, or for packaging. It is produced by pressing together moist fibers, typically cellulose pulp derived from wood, rags, or grasses, and drying them into flexible sheets.
Paperboard, sometimes known as cardboard, is generally thicker (usually over 0.25 mm or 10 points) than paper. According to ISO standards, paperboard is a paper with a basis weight (grammage) above 224 g/m2, but there are exceptions. Paperboard can be single- or multi-ply.
Corrugated fiberboard sometimes known as corrugated board or corrugated cardboard, is a combined paper-based material consisting of a fluted corrugated medium and one or two flat liner boards. The flute gives corrugated boxes much of their strength and is a contributing factor for why corrugated fiberboard is commonly used for shipping and storage.
There are also multiple names for containers:
A shipping container made of corrugated fiberboard is sometimes called a "cardboard box", a "carton", or a "case". There are many options for corrugated box design. Shipping container is used in shipping and transporting goods due to its strength and durability, thus corrugated boxes are designed to withstand the rigors of transportation and handling.
A folding carton made of paperboard is sometimes called a "cardboard box". Commonly used for packaging consumer goods, such as cereals, cosmetics, and pharmaceuticals. These cartons are designed to fold flat when empty, saving space during storage and transport.
A set-up box is made of a non-bending grade of paperboard and is sometimes called a "cardboard box". Often used for high-end products, such as jewelry, electronics, or gift items. Unlike folding cartons, set-up boxes do not fold flat and are delivered fully constructed.
Drink boxes made of paperboard laminates, are sometimes called "cardboard boxes", "cartons", or "boxes". Widely used for packaging beverages like juice, milk, and wine. These cartons are designed to maintain the freshness of liquid products and are often used in aseptic packaging.
History
The first commercial paperboard (not corrugated) box is sometimes credited to the firm M. Treverton & Son in England in 1817.
Cardboard box packaging was made the same year in Germany.
The Scottish-born Robert Gair invented the pre-cut cardboard or paperboard box in 1890 – flat pieces manufactured in bulk that folded into boxes. Gair's invention came about as a result of an accident: he was a Brooklyn printer and paper-bag maker during the 1870s, and one day, while he was printing an order of seed bags, a metal ruler normally used to crease bags shifted in position and cut them. Gair discovered that by cutting and creasing in one operation he could make prefabricated paperboard boxes. Applying this idea to corrugated boxboard was a straightforward development when the material became available around the turn of the twentieth century.
Cardboard boxes were developed in France about 1840 for transporting the Bombyx mori moth and its eggs by silk manufacturers, and for more than a century the manufacture of cardboard boxes was a major industry in the Valréas area.
The advent of lightweight flaked cereals increased the use of cardboard boxes. The first to use cardboard boxes as cereal cartons was the Kellogg Company.
Corrugated (also called pleated) paper was patented in England in 1856, and used as a liner for tall hats, but corrugated boxboard was not patented and used as a shipping material until 20 December 1871. The patent was issued to Albert Jones of New York City for single-sided (single-face) corrugated board.
Jones used the corrugated board for wrapping bottles and glass lantern chimneys. The first machine for producing large quantities of corrugated board was built in 1874 by G. Smyth, and in the same year Oliver Long improved upon Jones's design by inventing corrugated board with liner sheets on both sides.
This was corrugated cardboard as we know it today.
The first corrugated cardboard box manufactured in the US was in 1895.
By the early 1900s, wooden crates and boxes were being replaced by corrugated paper shipping cartons.
By 1908, the terms "corrugated paper-board" and "corrugated cardboard" were both in use in the paper trade.
Crafts and entertainment
Cardboard and other paper-based materials (paperboard, corrugated fiberboard, etc.) can have a post-primary life as a cheap material for the construction of a range of projects, among them being science experiments, children's toys, costumes, or insulative lining. Some children enjoy playing inside boxes.
A common cliché is that, if presented with a large and expensive new toy, a child will quickly become bored with the toy and play with the box instead. Although this is usually said somewhat jokingly, children certainly enjoy playing with boxes, using their imagination to portray the box as an infinite variety of objects. One example of this in popular culture is from the comic strip Calvin and Hobbes, whose protagonist, Calvin, often imagined a cardboard box as a "transmogrifier", a "duplicator", or a time machine.
So prevalent is the cardboard box's reputation as a plaything that in 2005 a cardboard box was added to the National Toy Hall of Fame in the US, one of very few non-brand-specific toys to be honoured with inclusion. As a result, a toy "house" (actually a log cabin) made from a large cardboard box was added to the Hall, housed at the Strong National Museum of Play in Rochester, New York.
The Metal Gear series of stealth video games has a running gag involving a cardboard box as an in-game item, which can be used by the player to try to sneak through places without getting caught by enemy sentries.
Housing and furniture
Living in a cardboard box is stereotypically associated with homelessness. However, in 2005, Melbourne architect Peter Ryan designed a house composed largely of cardboard.
More common are small seatings or little tables made from corrugated cardboard. Merchandise displays made of cardboard are often found in self-service shops.
Cushioning by crushing
Mass and viscosity of the enclosed air help together with the limited stiffness of boxes to absorb the energy of oncoming objects. In 2012, British stuntman Gary Connery safely landed via wingsuit without deploying his parachute, landing on a high crushable "runway" (landing zone) built with thousands of cardboard boxes.
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75820 | https://en.wikipedia.org/wiki/Cork%20%28material%29 | Cork (material) | Cork is an impermeable buoyant material. It is the phellem layer of bark tissue which is harvested for commercial use primarily from Quercus suber (the cork oak), which is native to southwest Europe and northwest Africa. Cork is composed of suberin, a hydrophobic substance. Because of its impermeable, buoyant, elastic, and fire retardant properties, it is used in a variety of products, the most common of which is wine stoppers.
The montado landscape of Portugal produces approximately half of the cork harvested annually worldwide, with Corticeira Amorim being the leading company in the industry. Cork was examined microscopically by Robert Hooke, which led to his discovery and naming of the cell.
Cork composition varies depending on geographic origin, climate and soil conditions, genetic origin, tree dimensions, age (virgin or reproduction), and growth conditions. However, in general, cork is made up of suberin (average of about 40%), lignin (22%), polysaccharides (cellulose and hemicellulose) (18%), extractables (15%) and others.
History
Cork is a natural material used by humans for over 5,000 years. It is a material whose applications have been known since antiquity, especially in floating devices and as stopper for beverages, mainly wine, whose market, from the early twentieth century, had a massive expansion, particularly due to the development of several cork-based agglomerates.
In China, Egypt, Babylon, and Persia from about 3000 BC, cork was already used for sealing containers, fishing equipment, and domestic applications. In ancient Greece (1600 to 1100 years BC) cork was used in footwear, to manufacture a type of sandals attached to the foot by straps, generally leather and with a sole in cork or leather.
In the second century AD, a Greek physician, Dioscorides, noted several medical applications of cork, mainly for hair loss treatment. Nowadays, the majority of people know cork for its use as stoppers in wine bottles. The innovation of using cork as stopper can be traced back to the late 17th century, attributed to Dom Pierre Pérignon. Cork stoppers were adopted in 1729 by Ruinart and in 1973 by Moët et Chandon.
Structure
Cork presents a characteristic cellular structure in which the cells have usually a pentagonal or hexagonal shape. The cellular wall consists of a thin, lignin-rich middle lamella (internal primary wall), a thick secondary wall made up from alternating suberin and wax lamella, and a thin tertiary wall of polysaccharides. Some studies suggest that the secondary wall is lignified, and therefore, may not consist exclusively of suberin and waxes. The cells of cork are filled with a gas mixture similar to air, making them behave as authentic "pads," which contributes to the capability of cork to recover after compression.
Sources
There are about 2,200,000 hectares of cork oak (Quercus suber) forest in the Mediterranean basin, the native area of the species. The most extensively managed habitats are in Portugal (34%) and in Spain (27%).
Annual production is about 300,000 tons; 49.6% from Portugal, 30.5% from Spain, 5.8% from Morocco, 4.9% from Algeria, 3.5% from Tunisia, 3.1% from Italy, and 2.6% from France.
Once the trees are about 25 years old the cork is traditionally stripped from the trunks every nine years, with the first two harvests generally producing lower quality cork (male cork or virgin cork). The trees live for about 300 years.
The cork industry is generally regarded as environmentally friendly. Cork production is generally considered sustainable because the cork tree is not cut down to obtain cork; only the bark is stripped to harvest the cork. The tree continues to live and grow. The sustainability of production and the easy recycling of cork products and by-products are two of its most distinctive aspects. Cork oak forests also prevent desertification and are a particular habitat in the Iberian Peninsula and the refuge of various endangered species.
Carbon footprint studies conducted by Corticeira Amorim, Oeneo Bouchage of France and the Cork Supply Group of Portugal concluded that cork is the most environmentally friendly wine stopper in comparison to other alternatives. The Corticeira Amorim's study, in particular ("Analysis of the life cycle of Cork, Aluminum and Plastic Wine Closures"), was developed by PricewaterhouseCoopers, according to ISO 14040. Results concluded that, concerning the emission of greenhouse gases, each plastic stopper released 10 times more CO2, whilst an aluminium screw cap releases 26 times more CO2 than does a cork stopper. For example, to produce 1,000 cork stoppers 1.5 kg are emitted, but to produce the same amount of plastic stoppers 14 kg of are emitted and for the same amount of aluminium screw caps 37 kg are emitted.
The Chinese cork oak is native to East Asia and is cultivated in a limited extent in China; the cork produced is considered inferior to Q. suber and are used to produce agglomerated cork products.
The so-called "cork trees" (Phellodendron) are unrelated to the cork oak, they have corky bark but not thick enough for cork production.
Harvesting
Cork is extracted only from early May to late August, when the cork can be separated from the tree without causing permanent damage. When the tree reaches 25–30 years of age and about 24 in (60 cm) in circumference, the cork can be removed for the first time. However, this first harvest almost always produces poor quality or virgin cork (Portuguese ; Spanish or ).
The workers who specialize in removing the cork are known as extractors. An extractor uses a very sharp axe to make two types of cuts on the tree: one horizontal cut around the plant, called a crown or necklace, at a height of about two to three times the circumference of the tree, and several vertical cuts called rulers or openings. This is the most delicate phase of the work because, even though cutting the cork requires significant force, the extractor must not damage the underlying phellogen or the tree will be harmed.
To free the cork from the tree, the extractor pushes the handle of the axe into the rulers. A good extractor needs to use a firm but precise touch in order to free a large amount of cork without damaging the product or tree.
These freed portions of the cork are called planks. The planks are usually carried off by hand since cork forests are rarely accessible to vehicles. The cork is stacked in piles in the forest or in yards at a factory and traditionally left to dry, after which it can be loaded onto a truck and shipped to a processor.
Bark from initial harvests can be used to make flooring, shoes, insulation and other industrial products. Subsequent extractions usually occur at intervals of nine years, though it can take up to thirteen for the cork to reach an acceptable size. If the product is of high quality it is known as gentle cork (Portuguese , but also only if it is the second time; Spanish , also restricted to the second time), and, ideally, is used to make stoppers for wine and champagne bottles.
Properties and uses
Cork's elasticity combined with its near-impermeability makes it suitable as a material for bottle stoppers, especially for wine bottles. Cork stoppers represent about 60% of all cork based production. Cork has an almost zero Poisson's ratio, which means the radius of a cork does not change significantly when squeezed or pulled.
Cork is an excellent gasket material. Some carburetor float bowl gaskets are made of cork, for example.
Cork is also an essential element in the production of badminton shuttlecocks.
Cork's bubble-form structure and natural fire retardant make it suitable for acoustic and thermal insulation in house walls, floors, ceilings, and facades. The by-product of more lucrative stopper production, corkboard, is gaining popularity as a non-allergenic, easy-to-handle and safe alternative to petrochemical-based insulation products.
Cork is also used to make vinyl record slipmats, due to its ability to not attract dust. They also dampen static and vibrations.
Sheets of cork, also often the by-product of stopper production, are used to make bulletin boards as well as floor and wall tiles.
Cork's low density makes it a suitable material for fishing floats and buoys, as well as handles for fishing rods (as an alternative to neoprene).
Granules of cork can also be mixed into concrete. The composites made by mixing cork granules and cement have lower thermal conductivity, lower density, and good energy absorption. Some of the property ranges of the composites are density (400–1500 kg/m3), compressive strength (1–26 MPa), and flexural strength (0.5–4.0 MPa).
Use in wine bottling
As late as the mid-17th century, French vintners did not use cork stoppers, using instead oil-soaked rags stuffed into the necks of bottles.
Wine corks can be made of either a single piece of cork, or composed of particles, as in champagne corks; corks made of granular particles are called "agglomerated corks".
Natural cork closures are used for about 80% of the 20 billion bottles of wine produced each year. After a decline in use as wine-stoppers due to the increase in the use of synthetic alternatives, cork wine-stoppers are making a comeback and currently represent approximately 60% of wine-stoppers in 2016.
Because of the cellular structure of cork, it is easily compressed upon insertion into a bottle and will expand to form a tight seal. The interior diameter of the neck of glass bottles tends to be inconsistent, making this ability to seal through variable contraction and expansion an important attribute. However, unavoidable natural flaws, channels, and cracks in the bark make the cork itself highly inconsistent. In a 2005 closure study, 45% of corks showed gas leakage during pressure testing both from the sides of the cork as well as through the cork body itself.
Since the mid-1990s, a number of wine brands have switched to alternative wine closures such as plastic stoppers, screw caps, or other closures. During 1972 more than half of the Australian bottled wine went bad due to corking. A great deal of anger and suspicion was directed at Portuguese and Spanish cork suppliers who were suspected of deliberately supplying bad cork to non-EEC wine makers to help prevent cheap imports. Cheaper wine makers developed the aluminium "Stelvin" cap with a polypropylene stopper wad. More expensive wines and carbonated varieties continued to use cork, although much closer attention was paid to the quality. Even so, some high premium makers prefer the Stelvin as it is a guarantee that the wine will be good even after many decades of ageing. Some consumers may have conceptions about screw caps being representative of lower quality wines, due to their cheaper price; however, in Australia, for example, much of the non-sparkling wine production now uses these Stelvin caps as a cork alternative, although some have recently switched back to cork citing issues using screw caps.
The alternatives to cork have both advantages and disadvantages. For example, screwtops are generally considered to offer a trichloroanisole (TCA) free seal, but they also reduce the oxygen transfer rate between the bottle and the atmosphere to almost zero, which can lead to a reduction in the quality of the wine. TCA is the main documented cause of cork taint in wine. However, some in the wine industry say natural cork stoppers are important because they allow oxygen to interact with wine for proper aging, and are best suited for wines purchased with the intent to age.
Stoppers which resemble natural cork very closely can be made by isolating the suberin component of the cork from the undesirable lignin, mixing it with the same substance used for contact lenses and an adhesive, and molding it into a standardized product, free of TCA or other undesirable substances. Composite corks with real cork veneers are used in cheaper wines. Celebrated Australian wine writer and critic James Halliday has written that since a cork placed inside the neck of a wine bottle is 350-year-old technology, it is logical to explore other more modern and precise methods of keeping wine safe.
The study "Analysis of the life cycle of Cork, Aluminum and Plastic Wine Closures," conducted by PricewaterhouseCoopers and commissioned by a major cork manufacturer, Amorim, concluded that cork is the most environmentally responsible stopper, in a one-year life cycle analysis comparison with plastic stoppers and aluminum screw caps.
Other uses
On 28 November 2007, the Portuguese national postal service CTT issued the world's first postage stamp made of cork.
In musical instruments, particularly woodwind instruments, where it is used to fasten together segments of the instrument, making the seams airtight. Low quality conducting baton handles are also often made out of cork.
In shoes, especially those using welt construction to improve climate control and comfort.
Because it is impermeable and moisture-resistant, cork is often used as an alternative to leather in handbags, wallets, and other fashion items.
To make bricks for the outer walls of houses, as in Portugal's pavilion at Expo 2000.
As the core of both baseballs and cricket balls. A corked bat is made by replacing the interior of a baseball bat with cork – a practice known as "corking". It was historically a method of cheating at baseball; the efficacy of the practice is now discredited.
In various forms, in spacecraft heat shields and fairings.
In the paper pick-up mechanisms in inkjet and laser printers.
To make later-model pith helmets.
Hung from hats to keep insects away. (See cork hat)
As a core material in sandwich composite construction.
As the friction lining material of an automatic transmission clutch, as designed in certain mopeds.
Alternative of wood or aluminium in automotive interiors.
Cork slabs are sometimes used by orchid growers as a natural mounting material.
Cork paddles are used by glass blowers to manipulate and shape hot molten glass.
Many racing bicycles have their handlebars wrapped in cork-based tape manufactured in a variety of colors.
To make architectural models.
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75838 | https://en.wikipedia.org/wiki/Pen | Pen | A pen is a common writing instrument that applies ink to a surface, usually paper, for writing or drawing. Early pens such as reed pens, quill pens, dip pens and ruling pens held a small amount of ink on a nib or in a small void or cavity that had to be periodically recharged by dipping the tip of the pen into an inkwell. Today, such pens find only a small number of specialized uses, such as in illustration and calligraphy. Reed pens, quill pens and dip pens, which were used for writing, have been replaced by ballpoint pens, rollerball pens, fountain pens and felt or ceramic tip pens. Ruling pens, which were used for technical drawing and cartography, have been replaced by technical pens such as the Rapidograph. All of these modern pens contain internal ink reservoirs, such that they do not need to be dipped in ink while writing.
Types
Modern
Pens commonly used today can be categorized based on the mechanism of the writing tip and the type of ink:
A ballpoint pen dispenses a viscous oil-based ink by means of a small hard sphere, or ball, which rolls over the surface being written on. The ball is held captive in a socket at the tip of the pen with one half exposed and the other half immersed in ink from the pen's reservoir. When the ball rotates, it transfers the ink - which wets the ball - from the reservoir to the external surface. The ball is typically under a millimeter in diameter and made of brass, steel, or tungsten carbide. The ink, due to its high viscosity, does not permeate through paper and does not leave the tip of the pen by capillary action. As such, a bare minimum amount of ink is dispensed, with the result that the writing dries almost instantly and ink lasts longer than it does in other types of pen. Ballpoint pens are reliable, versatile and robust, and are available for a very wide range of prices. They have replaced fountain pens as the most common tool for everyday writing.
A gel pen works similarly to a ballpoint pen, in that it dispenses ink using a rolling ball held in the writing tip. However, unlike oil-based ballpoint pen ink, gel pen ink consists of a water-based gel that has a pigment suspended in it. Because the ink is thick and opaque, it shows up more clearly on dark or slick surfaces than the typical inks used in ballpoint or felt tip pens. Gel pens can be used for many types of writing and illustration. Since the gel medium eliminates the constraints of a soluble dye, many new colors are made possible, as well as some special types of ink; gel pens are available in a wide range of vibrant or saturated colors, in pastel colors, in neon colors, in metallic colors, in glitter inks, in glow-in-the-dark ink, and so on.
A rollerball pen is a pen that dispenses a water-based ink through a ball tip similar to that of a ballpoint pen. As such, gel pens might be considered a subcategory of rollerball pens; however, due to the widespread knowledge and use of the term 'gel pen', 'rollerball' is in practice typically reserved for pens which use liquid ink. The lower viscosity of rollerball ink compared to oil-based ballpoint pen ink has several effects on the pen's performance. Since the ink flows more easily and is more easily absorbed into paper, more ink is dispensed in general. This changes the writing experience by lubricating the motion of the tip over the paper. It also results in a solid and uninterrupted line, since the diffusion of the ink through the paper fills small gaps that might otherwise be left by the ball point. Compared to ballpoint pens, which dispense a smaller amount of more viscous ink, the writing by a rollerball pen takes longer to dry on the page and can seep through thin paper such as to become visible on the opposite side. When the tip of a rollerball pen is held against paper, ink leaves the tip continually by capillary action in much the same way as would occur with a fountain pen. This can lead to ink blots or smears. The rollerball pen was initially designed to combine the convenience of a ballpoint pen with the smooth "wet ink" effect of a fountain pen. Refillable rollerball pens have recently become available; these generally use cartridges of fountain pen ink.
A felt-tip pen, or marker, has a porous tip made of fibrous material, which normally remains saturated with ink from the reservoir. As ink leaves the tip, new ink is drawn from the reservoir - which often consists of a large volume of a similar porous material to that used in the tip - by capillary action and gravity. As with a fountain pen, ink leaves the tip of a felt tip pen by capillary action when writing on a porous surface. However, unlike fountain pens, many markers can also reliably write on slick impermeable surfaces that are wet by the ink, and in such applications ink typically does not continually leave the pen as it is held against the writing surface. The smallest, finest-tipped felt-tip pens are used for writing on paper. Medium-sized felt-tips are often used by children for coloring and drawing. Larger types, often called "markers", are used for writing in larger sizes, often on surfaces other than paper such as corrugated boxes and whiteboards. Specialized felt-tip pens referred to by names such as "liquid chalk" or "chalkboard markers" are used to write on chalkboards. Markers with wide tips and bright but transparent ink, called highlighters, are used to highlight text that has already been written or printed. Pens designed for children or for temporary writing (as with a whiteboard or overhead projector) typically use non-permanent inks. Large markers used to label shipping cases or other packages are usually permanent markers.
A brush pen is a pen whose writing tip consists of a small brush fed with ink from a liquid ink reservoir similar to those used in fountain pens and rollerball pens. Brush pens might be either refillable or disposable, and might use either water-based or waterproof ink. The most significant functional difference of brush pens from felt-tip pens is the far greater compliance of the tip. Brush pens are an obvious alternative to ink brushes for Chinese calligraphy and Japanese calligraphy, but are now also commonly used in other forms of calligraphy and by artists such as illustrators and cartoonists. The primary appeal of these pens to such artists is that they allow a great deal of line width variation in response to small changes in applied pressure.
A stylus pen, plural styli or styluses, is a writing utensil which does not use ink, but rather makes marks primarily by creating scratches or indentations in the writing surface. As such, the tip often consists simply of a sharp metal point. Such tools are also used for other types of marking than writing, and for shaping or carving in, for example, pottery. The word stylus also refers to a pen-shaped computer accessory that is used to achieve greater precision when using touchscreens than generally possible with a fingertip. There are products available that combine a ballpoint tip at one end and a touchscreen stylus at the other.
Historic
These historic types of pens are no longer in common use as writing instruments, but may be used by calligraphers and other artists:
A fountain pen uses water-based liquid ink delivered through a nib, which is in general a flat piece of metal with a thin slit extending inwards from the writing tip. Driven by gravity, the ink flows from a reservoir to the nib through a feed, which is in general a specially shaped solid block of material with channels and grooves cut into it. The feed delivers the ink to the slit in the nib. While writing, ink is pulled out of this slit by capillary action. A fountain pen nib, unlike the tip of a ballpoint, gel or rollerball pen, has no moving parts. A fountain pen reservoir can be refillable or disposable; the disposable type is called an ink cartridge. A pen with a refillable reservoir may have a mechanism such as a piston to draw ink from a bottle through the nib, or it may require refilling with an eye dropper. Refillable reservoirs, also known as cartridge converters, are available for some pens otherwise designed to use disposable cartridges. A fountain pen can be used with permanent or non-permanent inks.
A dip pen (or nib pen) consists of a metal nib with capillary channels, like that of a fountain pen, mounted on a handle or holder, often made of wood. A dip pen is called such because it usually has no ink reservoir and must therefore be repeatedly dipped into an inkpot in order to recharge the nib with ink while drawing or writing. The dip pen has certain advantages over a fountain pen; it can use waterproof pigmented (particle-and-binder-based) inks, such as so-called India ink, drawing ink, or acrylic inks, which would destroy a fountain pen by clogging, as well as the traditional iron gall ink, which can cause corrosion in a fountain pen. Dip pens are now mainly used in illustration, calligraphy, and comics. A particularly fine-pointed type of dip pen known as a crowquill is a favorite instrument of artists such as David Stone Martin and Jay Lynch, because its flexible metal point can create a variety of delicate lines, textures and tones in response to variation of pressure while drawing.
The ink brush is the traditional writing implement in East Asian calligraphy. The body of the brush can be made from bamboo, or from rarer materials such as red sandalwood, glass, ivory, silver, and gold. The head of the brush can be made from the hair (or feathers) of a wide variety of animals, including the weasel, rabbit, deer, chicken, duck, goat, pig, and tiger. There is also a tradition both in China and in Japan of making a brush using the hair of a newborn, as a once-in-a-lifetime souvenir for the child. This practice is associated with the legend of an ancient Chinese scholar who ranked first in the Imperial examinations using such a personalized brush. Calligraphy brushes are widely considered an extension of the calligrapher's arm. Today, calligraphy may also be done using a pen, but pen calligraphy does not enjoy the same prestige as traditional brush calligraphy.
A quill is a pen made from a flight feather of a large bird, most often a goose. To make a quill, a feather must be cured through aging or heat-treatment, after which a nib is fashioned from the shaft by cutting a slit in it and carving away the sides to create a pointed tip. With practice, suitable feathers can be made into quills quickly and cheaply using no more than a small knife and a source of heat. Due to their easy availability, quills remained the writing instruments of choice in the West for a long time—from the 6th century to the 19th—before the metal dip pen, the fountain pen, and eventually the ballpoint pen came to be manufactured in large numbers. Quills, like later metal-nibbed dip pens, must periodically be dipped in ink while writing.
A reed pen is cut from a reed or bamboo, with a slit in a narrow tip. Its mechanism is essentially the same as that of a quill or a metal dip pen. The reed pen has almost disappeared but is still used by young school students in some parts of India and Pakistan, who learn to write with them on small timber boards known as "Takhti".
A paper pen, invented by Nasima Akhtar in 2007 from Jashore, Bangladesh, is an eco-friendly writing instrument made from paper. These pens are biodegradable and contain seeds at their base, allowing them to be planted after use to grow into various plants.
History
Ancient Egyptians had developed writing on papyrus scrolls when scribes used thin reed brushes or reed pens from the Juncus maritimus or sea rush. In his book A History of Writing, Steven Roger Fischer suggests, on the basis of finds at Saqqara, that the reed pen might well have been used for writing on parchment as long ago as the First Dynasty, or around 3000 BC. Reed pens continued to be used until the Middle Ages, but were slowly replaced by quills from about the 7th century. The reed pen, made from reed or bamboo, is still used in some parts of Pakistan by young students and is used to write on small wooden boards.
The reed pen survived until papyrus was replaced as a writing surface by animal skins, vellum and parchment. The smoother surface of skin allowed finer, smaller writing with a quill pen, derived from the flight feather. The quill pen was used in Qumran, Judea to write some of the Dead Sea Scrolls, which date back to around 100 BC. The scrolls were written in Hebrew dialects with bird feathers or quills. There is a specific reference to quills in the writings of St. Isidore of Seville in the 7th century. Quill pens were still widely used in the eighteenth century, and were used to write and sign the Constitution of the United States in 1787.
A copper nib was found in the ruins of Pompeii, showing that metal nibs were used in the year 79. There is also a reference to 'a silver pen to carry ink in', in Samuel Pepys' diary for August 1663. 'New invented' metal pens are advertised in The Times in 1792. A metal pen point was patented in 1803, but the patent was not commercially exploited. A patent for the manufacture of metal pens was advertised for sale by Bryan Donkin in 1811. John Mitchell of Birmingham started to mass-produce pens with metal nibs in 1822, and after that, the quality of steel nibs improved enough so that dip pens with metal nibs came into general use.
The earliest historical record of a pen with a reservoir dates back to the 10th century AD. In 953, Ma'ād al-Mu'izz, the Fatimid Caliph of Egypt, demanded a pen which would not stain his hands or clothes, and was provided with a pen which held ink in a reservoir and delivered it to the nib. This pen may have been a fountain pen, but its mechanism remains unknown, and only one record mentioning it has been found. A later reservoir pen was developed in 1636. In his Deliciae Physico-Mathematicae (1636), German inventor Daniel Schwenter described a pen made from two quills. One quill served as a reservoir for ink inside the other quill. The ink was sealed inside the quill with cork. Ink was squeezed through a small hole to the writing point. In 1809, Bartholomew Folsch received a patent in England for a pen with an ink reservoir.
A student in Paris, Romanian Petrache Poenaru invented a fountain pen that used a quill as an ink reservoir. The French Government patented this in May 1827. Fountain pen patents and production then increased in the 1850s.
The first patent on a ballpoint pen was issued on October 30, 1888, to John J. Loud. In 1938, László Bíró, a Hungarian newspaper editor, with the help of his brother George, a chemist, began to design new types of pens, including one with a tiny ball in its tip that was free to turn in a socket. As the pen moved along the paper, the ball rotated, picking up ink from the ink cartridge and leaving it on the paper. Bíró filed a British patent on June 15, 1938. In 1940, the Bíró brothers and a friend, Juan Jorge Meyne, moved to Argentina, fleeing Nazi Germany. On June 17, 1943, they filed for another patent. They formed "Bíró Pens of Argentina", and by the summer of 1943, the first commercial models were available. Erasable ballpoint pens were introduced by Paper Mate in 1979, when the Erasermate was put on the market.
Slavoljub Eduard Penkala, a Croatian engineer and inventor, became renowned for further development of the mechanical pencil (1906) – then called an "automatic pencil" – and the first solid-ink fountain pen (1907). Collaborating with the Croatian entrepreneur Edmund Moster, he started the Penkala-Moster Company and built a pen-and-pencil factory that was one of the biggest in the world at the time. This company, now called TOZ-Penkala, still exists today. "TOZ" stands for "", meaning "Zagreb Pencil Factory".
In the 1960s, the fiber- or felt-tipped pen was invented by Yukio Horie of the Tokyo Stationery Company, Japan. Paper Mate's Flair was among the first felt-tip pens to hit the U.S. market in the 1960s, and it has been the leader ever since. Marker pens and highlighters, both similar to felt pens, have become popular in recent times.
Rollerball pens were introduced in the early 1970s. They use a mobile ball and liquid ink to produce a smoother line. Technological advances during the late 1980s and early 1990s have improved the roller ball's overall performance. A porous point pen contains a point made of some porous material such as felt or ceramic. A high quality drafting pen will usually have a ceramic tip, since this wears well and does not broaden when pressure is applied while writing.
Although the invention of the typewriter and personal computer with the keyboard input method has offered another way to write, the pen is still the main means of writing. Many people like to use expensive types and brands of pens, including fountain pens, and these are sometimes regarded as a status symbol.
| Technology | Artist's tools | null |
75873 | https://en.wikipedia.org/wiki/Nernst%20equation | Nernst equation | In electrochemistry, the Nernst equation is a chemical thermodynamical relationship that permits the calculation of the reduction potential of a reaction (half-cell or full cell reaction) from the standard electrode potential, absolute temperature, the number of electrons involved in the redox reaction, and activities (often approximated by concentrations) of the chemical species undergoing reduction and oxidation respectively. It was named after Walther Nernst, a German physical chemist who formulated the equation.
Expression
General form with chemical activities
When an oxidizer () accepts a number z of electrons () to be converted in its reduced form (), the half-reaction is expressed as:
Ox + ze- -> Red
The reaction quotient (), also often called the ion activity product (IAP), is the ratio between the chemical activities (a) of the reduced form (the reductant, ) and the oxidized form (the oxidant, ). The chemical activity of a dissolved species corresponds to its true thermodynamic concentration taking into account the electrical interactions between all ions present in solution at elevated concentrations. For a given dissolved species, its chemical activity (a) is the product of its activity coefficient (γ) by its molar (mol/L solution), or molal (mol/kg water), concentration (C): a = γ C. So, if the concentration (C, also denoted here below with square brackets [ ]) of all the dissolved species of interest are sufficiently low and that their activity coefficients are close to unity, their chemical activities can be approximated by their concentrations as commonly done when simplifying, or idealizing, a reaction for didactic purposes:
At chemical equilibrium, the ratio of the activity of the reaction product (aRed) by the reagent activity (aOx) is equal to the equilibrium constant of the half-reaction:
The standard thermodynamics also says that the actual Gibbs free energy is related to the free energy change under standard state by the relationship:
where is the reaction quotient and R is the universal ideal gas constant.
The cell potential associated with the electrochemical reaction is defined as the decrease in Gibbs free energy per coulomb of charge transferred, which leads to the relationship The constant (the Faraday constant) is a unit conversion factor , where is the Avogadro constant and is the fundamental electron charge. This immediately leads to the Nernst equation, which for an electrochemical half-cell is
For a complete electrochemical reaction (full cell), the equation can be written as
where:
is the half-cell reduction potential at the temperature of interest,
is the standard half-cell reduction potential,
is the cell potential (electromotive force) at the temperature of interest,
is the standard cell potential,
is the universal ideal gas constant: ,
is the temperature in kelvins,
is the number of electrons transferred in the cell reaction or half-reaction,
is the Faraday constant, the magnitude of charge (in coulombs) per mole of electrons: ,
is the reaction quotient of the cell reaction, and,
is the chemical activity for the relevant species, where is the activity of the reduced form and is the activity of the oxidized form.
Thermal voltage
At room temperature (25 °C), the thermal voltage is approximately 25.693 mV. The Nernst equation is frequently expressed in terms of base-10 logarithms (i.e., common logarithms) rather than natural logarithms, in which case it is written:
where λ = ln(10) ≈ 2.3026 and λVT ≈ 0.05916 Volt.
Form with activity coefficients and concentrations
Similarly to equilibrium constants, activities are always measured with respect to the standard state (1 mol/L for solutes, 1 atm for gases, and T = 298.15 K, i.e., 25 °C or 77 °F). The chemical activity of a species , , is related to the measured concentration via the relationship , where is the activity coefficient of the species . Because activity coefficients tend to unity at low concentrations, or are unknown or difficult to determine at medium and high concentrations, activities in the Nernst equation are frequently replaced by simple concentrations and then, formal standard reduction potentials used.
Taking into account the activity coefficients () the Nernst equation becomes:
Where the first term including the activity coefficients () is denoted and called the formal standard reduction potential, so that can be directly expressed as a function of and the concentrations in the simplest form of the Nernst equation:
Formal standard reduction potential
When wishing to use simple concentrations in place of activities, but that the activity coefficients are far from unity and can no longer be neglected and are unknown or too difficult to determine, it can be convenient to introduce the notion of the "so-called" standard formal reduction potential () which is related to the standard reduction potential as follows:
So that the Nernst equation for the half-cell reaction can be correctly formally written in terms of concentrations as:
and likewise for the full cell expression.
According to Wenzel (2020), a formal reduction potential is the reduction potential that applies to a half reaction under a set of specified conditions such as, e.g., pH, ionic strength, or the concentration of complexing agents.
The formal reduction potential is often a more convenient, but conditional, form of the standard reduction potential, taking into account activity coefficients and specific conditions characteristics of the reaction medium. Therefore, its value is a conditional value, i.e., that it depends on the experimental conditions and because the ionic strength affects the activity coefficients, will vary from medium to medium. Several definitions of the formal reduction potential can be found in the literature, depending on the pursued objective and the experimental constraints imposed by the studied system. The general definition of refers to its value determined when . A more particular case is when is also determined at pH 7, as e.g. for redox reactions important in biochemistry or biological systems.
Determination of the formal standard reduction potential when 1
The formal standard reduction potential can be defined as the measured reduction potential of the half-reaction at unity concentration ratio of the oxidized and reduced species (i.e., when 1) under given conditions.
Indeed:
as, , when ,
, when ,
because , and that the term is included in .
The formal reduction potential makes possible to more simply work with molar (mol/L, M) or molal (mol/kg , m) concentrations in place of activities. Because molar and molal concentrations were once referred as formal concentrations, it could explain the origin of the adjective formal in the expression formal potential.
The formal potential is thus the reversible potential of an electrode at equilibrium immersed in a solution where reactants and products are at unit concentration. If any small incremental change of potential causes a change in the direction of the reaction, i.e. from reduction to oxidation or vice versa, the system is close to equilibrium, reversible and is at its formal potential. When the formal potential is measured under standard conditions (i.e. the activity of each dissolved species is 1 mol/L, T = 298.15 K = 25 °C = 77 °F, = 1 bar) it becomes de facto a standard potential. According to Brown and Swift (1949):
"A formal potential is defined as the potential of a half-cell, measured against the standard hydrogen electrode, when the total concentration of each oxidation state is one formal".
In this case, as for the standard reduction potentials, the concentrations of dissolved species remain equal to one molar (M) or one molal (m), and so are said to be one formal (F). So, expressing the concentration in molarity (1 mol/L):
The term formal concentration (F) is now largely ignored in the current literature and can be commonly assimilated to molar concentration (M), or molality (m) in case of thermodynamic calculations.
The formal potential is also found halfway between the two peaks in a cyclic voltammogram, where at this point the concentration of Ox (the oxidized species) and Red (the reduced species) at the electrode surface are equal.
The activity coefficients and are included in the formal potential , and because they depend on experimental conditions such as temperature, ionic strength, and pH, cannot be referred as an immutable standard potential but needs to be systematically determined for each specific set of experimental conditions.
Formal reduction potentials are applied to simplify calculations of a considered system under given conditions and measurements interpretation. The experimental conditions in which they are determined and their relationship to the standard reduction potentials must be clearly described to avoid to confuse them with standard reduction potentials.
Formal standard reduction potential at pH 7
Formal standard reduction potentials () are also commonly used in biochemistry and cell biology for referring to standard reduction potentials measured at pH 7, a value closer to the pH of most physiological and intracellular fluids than the standard state pH of 0. The advantage is to defining a more appropriate redox scale better corresponding to real conditions than the standard state. Formal standard reduction potentials () allow to more easily estimate if a redox reaction supposed to occur in a metabolic process or to fuel microbial activity under some conditions is feasible or not.
While, standard reduction potentials always refer to the standard hydrogen electrode (SHE), with [] = 1 M corresponding to a pH 0, and fixed arbitrarily to zero by convention, it is no longer the case at a pH of 7. Then, the reduction potential of a hydrogen electrode operating at pH 7 is -0.413 V with respect to the standard hydrogen electrode (SHE).
Expression of the Nernst equation as a function of pH
The and pH of a solution are related by the Nernst equation as commonly represented by a Pourbaix diagram . explicitly denotes expressed versus the standard hydrogen electrode (SHE). For a half cell equation, conventionally written as a reduction reaction (i.e., electrons accepted by an oxidant on the left side):
The half-cell standard reduction potential is given by
where is the standard Gibbs free energy change, is the number of electrons involved, and is the Faraday's constant. The Nernst equation relates pH and as follows:
where curly brackets indicate activities, and exponents are shown in the conventional manner. This equation is the equation of a straight line for as a function of pH with a slope of volt (pH has no units).
This equation predicts lower at higher pH values. This is observed for the reduction of O2 into H2O, or OH−, and for the reduction of H+ into H2. is then often noted as to indicate that it refers to the standard hydrogen electrode (SHE) whose = 0 by convention under standard conditions (T = 298.15 K = 25 °C = 77 F, Pgas = 1 atm (1.013 bar), concentrations = 1 M and thus pH = 0).
Main factors affecting the formal standard reduction potentials
The main factor affecting the formal reduction potentials in biochemical or biological processes is most often the pH. To determine approximate values of formal reduction potentials, neglecting in a first approach changes in activity coefficients due to ionic strength, the Nernst equation has to be applied taking care to first express the relationship as a function of pH. The second factor to be considered are the values of the concentrations taken into account in the Nernst equation. To define a formal reduction potential for a biochemical reaction, the pH value, the concentrations values and the hypotheses made on the activity coefficients must always be explicitly indicated. When using, or comparing, several formal reduction potentials they must also be internally consistent.
Problems may occur when mixing different sources of data using different conventions or approximations (i.e., with different underlying hypotheses). When working at the frontier between inorganic and biological processes (e.g., when comparing abiotic and biotic processes in geochemistry when microbial activity could also be at work in the system), care must be taken not to inadvertently directly mix standard reduction potentials versus SHE (pH = 0) with formal reduction potentials (pH = 7). Definitions must be clearly expressed and carefully controlled, especially if the sources of data are different and arise from different fields (e.g., picking and mixing data from classical electrochemistry and microbiology textbooks without paying attention to the different conventions on which they are based).
Examples with a Pourbaix diagram
To illustrate the dependency of the reduction potential on pH, one can simply consider the two oxido-reduction equilibria determining the water stability domain in a Pourbaix diagram . When water is submitted to electrolysis by applying a sufficient difference of electrical potential between two electrodes immersed in water, hydrogen is produced at the cathode (reduction of water protons) while oxygen is formed at the anode (oxidation of water oxygen atoms). The same may occur if a reductant stronger than hydrogen (e.g., metallic Na) or an oxidant stronger than oxygen (e.g., F2) enters in contact with water and reacts with it. In the here beside (the simplest possible version of a Pourbaix diagram), the water stability domain (grey surface) is delimited in term of redox potential by two inclined red dashed lines:
Lower stability line with hydrogen gas evolution due to the proton reduction at very low Eh:
(cathode: reduction)
Higher stability line with oxygen gas evolution due to water oxygen oxidation at very high Eh:
(anode: oxidation)
When solving the Nernst equation for each corresponding reduction reaction (need to revert the water oxidation reaction producing oxygen), both equations have a similar form because the number of protons and the number of electrons involved within a reaction are the same and their ratio is one (2/2 for H2 and 4/4 with respectively), so it simplifies when solving the Nernst equation expressed as a function of pH.
The result can be numerically expressed as follows:
Note that the slopes of the two water stability domain upper and lower lines are the same (-59.16 mV/pH unit), so they are parallel on a Pourbaix diagram. As the slopes are negative, at high pH, both hydrogen and oxygen evolution requires a much lower reduction potential than at low pH.
For the reduction of H+ into H2 the here above mentioned relationship becomes:
because by convention = 0 V for the standard hydrogen electrode (SHE: pH = 0). So, at pH = 7, = -0.414 V for the reduction of protons.
For the reduction of O2 into 2 H2O the here above mentioned relationship becomes:
because = +1.229 V with respect to the standard hydrogen electrode (SHE: pH = 0). So, at pH = 7, = +0.815 V for the reduction of oxygen.
The offset of -414 mV in is the same for both reduction reactions because they share the same linear relationship as a function of pH and the slopes of their lines are the same. This can be directly verified on a Pourbaix diagram. For other reduction reactions, the value of the formal reduction potential at a pH of 7, commonly referred for biochemical reactions, also depends on the slope of the corresponding line in a Pourbaix diagram i.e. on the ratio of the number of to the number of involved in the reduction reaction, and thus on the stoichiometry of the half-reaction. The determination of the formal reduction potential at pH = 7 for a given biochemical half-reaction requires thus to calculate it with the corresponding Nernst equation as a function of pH. One cannot simply apply an offset of -414 mV to the Eh value (SHE) when the ratio differs from 1.
Applications in biology
Beside important redox reactions in biochemistry and microbiology, the Nernst equation is also used in physiology for calculating the electric potential of a cell membrane with respect to one type of ion. It can be linked to the acid dissociation constant.
Nernst potential
The Nernst equation has a physiological application when used to calculate the potential of an ion of charge across a membrane. This potential is determined using the concentration of the ion both inside and outside the cell:
When the membrane is in thermodynamic equilibrium (i.e., no net flux of ions), and if the cell is permeable to only one ion, then the membrane potential must be equal to the Nernst potential for that ion.
Goldman equation
When the membrane is permeable to more than one ion, as is inevitably the case, the resting potential can be determined from the Goldman equation, which is a solution of G-H-K influx equation under the constraints that total current density driven by electrochemical force is zero:
where
is the membrane potential (in volts, equivalent to joules per coulomb),
is the permeability for that ion (in meters per second),
is the extracellular concentration of that ion (in moles per cubic meter, to match the other SI units, though the units strictly don't matter, as the ion concentration terms become a dimensionless ratio),
is the intracellular concentration of that ion (in moles per cubic meter),
is the ideal gas constant (joules per kelvin per mole),
is the temperature in kelvins,
is the Faraday's constant (coulombs per mole).
The potential across the cell membrane that exactly opposes net diffusion of a particular ion through the membrane is called the Nernst potential for that ion. As seen above, the magnitude of the Nernst potential is determined by the ratio of the concentrations of that specific ion on the two sides of the membrane. The greater this ratio the greater the tendency for the ion to diffuse in one direction, and therefore the greater the Nernst potential required to prevent the diffusion. A similar expression exists that includes (the absolute value of the transport ratio). This takes transporters with unequal exchanges into account. See: sodium-potassium pump where the transport ratio would be 2/3, so r equals 1.5 in the formula below. The reason why we insert a factor r = 1.5 here is that current density by electrochemical force Je.c.(Na+) + Je.c.(K+) is no longer zero, but rather Je.c.(Na+) + 1.5Je.c.(K+) = 0 (as for both ions flux by electrochemical force is compensated by that by the pump, i.e. Je.c. = −Jpump), altering the constraints for applying GHK equation. The other variables are the same as above. The following example includes two ions: potassium (K+) and sodium (Na+). Chloride is assumed to be in equilibrium.
When chloride (Cl−) is taken into account,
Derivation
Using Boltzmann factor
For simplicity, we will consider a solution of redox-active molecules that undergo a one-electron reversible reaction
and that have a standard potential of zero, and in which the activities are well represented by the concentrations (i.e. unit activity coefficient). The chemical potential of this solution is the difference between the energy barriers for taking electrons from and for giving electrons to the working electrode that is setting the solution's electrochemical potential. The ratio of oxidized to reduced molecules, , is equivalent to the probability of being oxidized (giving electrons) over the probability of being reduced (taking electrons), which we can write in terms of the Boltzmann factor for these processes:
Taking the natural logarithm of both sides gives
If at = 1, we need to add in this additional constant:
Dividing the equation by to convert from chemical potentials to electrode potentials, and remembering that , we obtain the Nernst equation for the one-electron process :
Using thermodynamics (chemical potential)
Quantities here are given per molecule, not per mole, and so Boltzmann constant and the electron charge are used instead of the gas constant and Faraday's constant . To convert to the molar quantities given in most chemistry textbooks, it is simply necessary to multiply by the Avogadro constant: and . The entropy of a molecule is defined as
where is the number of states available to the molecule. The number of states must vary linearly with the volume of the system (here an idealized system is considered for better understanding, so that activities are posited very close to the true concentrations). Fundamental statistical proof of the mentioned linearity goes beyond the scope of this section, but to see this is true it is simpler to consider usual isothermal process for an ideal gas where the change of entropy takes place. It follows from the definition of entropy and from the condition of constant temperature and quantity of gas that the change in the number of states must be proportional to the relative change in volume . In this sense there is no difference in statistical properties of ideal gas atoms compared with the dissolved species of a solution with activity coefficients equaling one: particles freely "hang around" filling the provided volume), which is inversely proportional to the concentration , so we can also write the entropy as
The change in entropy from some state 1 to another state 2 is therefore
so that the entropy of state 2 is
If state 1 is at standard conditions, in which is unity (e.g., 1 atm or 1 M), it will merely cancel the units of . We can, therefore, write the entropy of an arbitrary molecule A as
where is the entropy at standard conditions and [A] denotes the concentration of A. The change in entropy for a reaction
is then given by
We define the ratio in the last term as the reaction quotient:
where the numerator is a product of reaction product activities, , each raised to the power of a stoichiometric coefficient, , and the denominator is a similar product of reactant activities. All activities refer to a time . Under certain circumstances (see chemical equilibrium) each activity term such as may be replaced by a concentration term, [A].In an electrochemical cell, the cell potential is the chemical potential available from redox reactions (). is related to the Gibbs free energy change only by a constant:
, where is the number of electrons transferred and is the Faraday constant. There is a negative sign because a spontaneous reaction has a negative Gibbs free energy and a positive potential . The Gibbs free energy is related to the entropy by , where is the enthalpy and is the temperature of the system. Using these relations, we can now write the change in Gibbs free energy,
and the cell potential,
This is the more general form of the Nernst equation.
For the redox reaction ,
and we have:
The cell potential at standard temperature and pressure (STP) is often replaced by the formal potential , which includes the activity coefficients of the dissolved species under given experimental conditions (T, P, ionic strength, pH, and complexing agents) and is the potential that is actually measured in an electrochemical cell.
Relation to the chemical equilibrium
The standard Gibbs free energy is related to the equilibrium constant as follows:
At the same time, is also equal to the product of the total charge () transferred during the reaction and the cell potential ():
The sign is negative, because the considered system performs the work and thus releases energy.
So,
And therefore:
Starting from the Nernst equation, one can also demonstrate the same relationship in the reverse way.
At chemical equilibrium, or thermodynamic equilibrium, the electrochemical potential and therefore the reaction quotient () attains the special value known as the equilibrium constant ():
Therefore,
Or at standard state,
We have thus related the standard electrode potential and the equilibrium constant of a redox reaction.
Limitations
In dilute solutions, the Nernst equation can be expressed directly in the terms of concentrations (since activity coefficients are close to unity). But at higher concentrations, the true activities of the ions must be used. This complicates the use of the Nernst equation, since estimation of non-ideal activities of ions generally requires experimental measurements. The Nernst equation also only applies when there is no net current flow through the electrode. The activity of ions at the electrode surface changes when there is current flow, and there are additional overpotential and resistive loss terms which contribute to the measured potential.
At very low concentrations of the potential-determining ions, the potential predicted by Nernst equation approaches toward . This is physically meaningless because, under such conditions, the exchange current density becomes very low, and there may be no thermodynamic equilibrium necessary for Nernst equation to hold. The electrode is called unpoised in such case. Other effects tend to take control of the electrochemical behavior of the system, like the involvement of the solvated electron in electricity transfer and electrode equilibria, as analyzed by Alexander Frumkin and B. Damaskin, Sergio Trasatti, etc.
Time dependence of the potential
The expression of time dependence has been established by Karaoglanoff.
Significance in other scientific fields
The Nernst equation has been involved in the scientific controversy about cold fusion. Fleischmann and Pons, claiming that cold fusion could exist, calculated that a palladium cathode immersed in a heavy water electrolysis cell could achieve up to 1027 atmospheres of pressure inside the crystal lattice of the metal of the cathode, enough pressure to cause spontaneous nuclear fusion. In reality, only 10,000–20,000 atmospheres were achieved. The American physicist John R. Huizenga claimed their original calculation was affected by a misinterpretation of the Nernst equation. He cited a paper about Pd–Zr alloys.
The Nernst equation allows the calculation of the extent of reaction between two redox systems and can be used, for example, to assess whether a particular reaction will go to completion or not. At chemical equilibrium, the electromotive forces (emf) of the two half cells are equal. This allows the equilibrium constant of the reaction to be calculated and hence the extent of the reaction.
| Physical sciences | Electrochemistry | Chemistry |
75887 | https://en.wikipedia.org/wiki/Webmail | Webmail | Webmail (or web-based email) is an email service that can be accessed using a standard web browser. It contrasts with email service accessible through a specialised email client software. Additionally, many internet service providers (ISP) provide webmail as part of their internet service package. Similarly, some web hosting providers also provide webmail as a part of their hosting package.
As with any web application, webmail's main advantage over the use of a desktop email client is the ability to send and receive email anywhere from a web browser.
History
Early implementations
The first Web Mail implementation was developed at CERN in 1993 by Phillip Hallam-Baker as a test of the HTTP protocol stack, but was not developed further. In the next two years, however, several people produced working webmail applications.
In Europe, there were three implementations, Søren Vejrum's "WWW Mail", Luca Manunza's "WebMail", and Remy Wetzels' "WebMail". Søren Vejrum's "WWW Mail" was written when he was studying and working at the Copenhagen Business School in Denmark, and was released on February 28, 1995. Luca Manunza's "WebMail" was written while he was working at CRS4 in Sardinia, from an idea of Gianluigi Zanetti, with the first source release on March 30, 1995. Remy Wetzels' "WebMail" was written while he was studying at the Eindhoven University of Technology in the Netherlands for the DSE and was released early January 1995.
In the United States, Matt Mankins wrote "Webex", and Bill Fitler, while at Lotus cc:Mail, began working on an implementation which he demonstrated publicly at Lotusphere on January 24, 1995. Matt Mankins, under the supervision of Dr. Burt Rosenberg at the University of Miami, released his "Webex" application source code in a post to comp.mail.misc on August 8, 1995, although it had been in use as the primary email application at the School of Architecture where Mankins worked for some months prior.
Bill Fitler's webmail implementation was further developed as a commercial product, which Lotus announced and released in the fall of 1995 as cc:Mail for the World Wide Web 1.0; thereby providing an alternative means of accessing a cc:Mail message store (the usual means being a cc:Mail desktop application that operated either via dialup or within the confines of a local area network).
Early commercialization of webmail was also achieved when "Webex" began to be sold by Mankins' company, DotShop, Inc., at the end of 1995. Within DotShop, "Webex" changed its name to "EMUmail"; which would be sold to companies like UPS and Rackspace until its sale to Accurev in 2001. EMUmail was one of the first applications to feature a free version that included embedded advertising, as well as a licensed version that did not.
Hotmail and Four11's RocketMail both launched in 1996 as free services and immediately became very popular.
Widespread deployment
As the 1990s progressed, and into the 2000s, it became more common for the general public to have access to webmail because:
many Internet service providers (such as EarthLink) and web hosting providers (such as Verio) began bundling webmail into their service offerings (often in parallel with POP/SMTP services);
many other enterprises (such as universities and large corporations) also started offering webmail as a way for their user communities to access their email (either locally managed or outsourced);
webmail service providers (such as Hotmail and RocketMail) emerged in 1996 as a free service to the general public, and rapidly gained in popularity.
In some cases, webmail application software is developed in-house by the organizations running and managing the application, and in some cases it is obtained from software companies that develop and sell such applications, usually as part of an integrated mail server package (an early example being Netscape Messaging Server). The market for webmail application software has continued into the 2010s.
Rendering and compatibility
Email users may find the use of both a webmail client and a desktop client using the POP3 protocol presents some difficulties. For example, email messages that are downloaded by the desktop client and are removed from the server will no longer be available on the webmail client. The user is limited to previewing messages using the web client before they are downloaded by the desktop email client. However, one may choose to leave the emails on the server, in which case this problem does not occur. The use of both a webmail client and a desktop client using the IMAP4 protocol allows the contents of the mailbox to be consistently displayed in both the webmail and desktop clients and any action the user performs on messages in one interface will be reflected when the email is accessed via the other interface. There are significant differences in rendering capabilities for many popular webmail services such as Gmail, Outlook.com and Yahoo! Mail. Due to the varying treatment of HTML tags, such as <style> and <head>, as well as CSS rendering inconsistencies, email marketing companies rely on older web development techniques to send cross-platform mail. This usually means a greater reliance on tables and inline stylesheets.
Microsoft Windows applications by default create email messages via MAPI. Several vendors produce tools to provide a MAPI interface to webmail.
Privacy concerns
Although emails stored unencrypted on any service provider's servers can be read by that service provider, specific concerns have been raised regarding webmail services that automatically analyze the contents of users' emails for the purpose of targeted advertising. At least two such services, Gmail and Yahoo! Mail, give users the option to opt out of targeted advertising.
Webmail that is accessed over unsecured HTTP may be readable by a third party who has access to the data transmission, such as over an unsecured Wi-Fi connection. This may be avoided by connecting to the webmail service via HTTPS, which encrypts the connection. Gmail has supported HTTPS since launch and in 2014 began requiring it for all webmail connections. Yahoo! Mail added the option to connect over HTTPS in 2013 and made HTTPS required in 2014.
| Technology | Computer software | null |
75971 | https://en.wikipedia.org/wiki/Sabre | Sabre | A sabre or saber ( ) is a type of backsword with a curved blade associated with the light cavalry of the early modern and Napoleonic periods. Originally associated with Central European cavalry such as the hussars, the sabre became widespread in Western Europe during the Thirty Years' War. Lighter sabres also became popular with infantry of the early 17th century. In the 19th century, models with less curving blades became common and were also used by heavy cavalry.
The military sabre was used as a duelling weapon in academic fencing in the 19th century, giving rise to a discipline of modern sabre fencing (introduced in the 1896 Summer Olympics) loosely based on the characteristics of the historical weapon, although in Olympic fencing, only cuts are allowed.
Etymology
The English sabre is recorded from the 1670s, as a direct loan from French, where sabre is an alteration of sable, which was in turn loaned from German Säbel, Sabel in the 1630s. The German word is on record from the 15th century, loaned from Polish szabla, which was itself adopted from Hungarian szabla (14th century, later szablya).
The spread of the Hungarian word to neighboring European languages took place in the context of the Ottoman wars in Europe of the 15th to 17th centuries. The spelling saber became common in American English in the second half of the 19th century.
The origin of the Hungarian word is unclear. It may itself be a loan from South Slavic (Serbo-Croatian sablja, Common Slavic *sabľa), which would ultimately derive from a Turkic source. In a more recent suggestion, the Hungarian word may ultimately derive from a Tungusic source, via Kipchak Turkic selebe, with later metathesis (of l-b to b-l) and apocope changed to *seble, which would have changed its vocalisation in Hungarian to the recorded sabla, perhaps under the influence of the Hungarian word szab- "to crop; cut (into shape)".
History
Origins
Though single-edged cutting swords already existed in the Ancient world, such as the ancient Egyptian and Sumerian sickle swords, these (usually forward instead of backward curving) weapons were chopping weapons for foot soldiers. This type of weapon developed into such heavy chopping weapons as the Greek Machaira and Anatolian Drepanon, and it still survives as the heavy Kukri chopping knife of the Gurkhas. However, in ancient China foot soldiers and cavalry often used a straight, single edged sword, and in the sixth century CE a longer, slightly curved cavalry variety of this weapon appeared in southern Siberia. This "proto-sabre" (the Turko-Mongol sabre) had developed into the true cavalry sabre by the eight century CE, and by the ninth century, it had become the usual side arm on the Eurasian steppes. The sabre arrived in Europe with the Magyars and the Turkic expansion. These oldest sabres had a slight curve, short, down-turned quillons, the grip facing the opposite direction to the blade and a sharp point with the top third of the reverse edge sharpened.
Early modern period
The introduction of the sabre proper in Western Europe, along with the term sabre itself, dates to the 17th century, via the influence of the szabla type ultimately derived from these medieval backswords. The adoption of the term is connected to the employment of Hungarian hussar (huszár) cavalry by Western European armies at the time. Hungarian hussars were employed as light cavalry, with the role of harassing enemy skirmishers, overrunning artillery positions, and pursuing fleeing troops. In the late 17th and early 18th centuries, many Hungarian hussars fled to other Central and Western European countries and became the core of light cavalry formations created there. The Hungarian term szablya is ultimately traced to the northwestern Turkic selebe, with contamination from the Hungarian verb szab "to cut".
The original type of sabre, or Polish szabla, was used as a cavalry weapon, possibly inspired by Hungarian or wider Turco-Mongol warfare.
The karabela was a type of szabla popular in the late 17th century, worn by the Polish–Lithuanian Commonwealth nobility class, the szlachta. While designed as a cavalry weapon, it also came to replace various types of straight-bladed swords used by infantry. The Swiss sabre originated as a regular sword with a single-edged blade in the early 16th century, but by the 17th century began to exhibit specialized hilt types.
Polish–Lithuanian Commonwealth
In the Polish–Lithuanian Commonwealth (16th–18th century) a specific type of sabre-like melee weapon, the szabla, was used. Richly decorated sabres were popular among the Polish nobility, who considered it to be one of the most important pieces of men's traditional attire. With time, the design of the sabre greatly evolved in the commonwealth and gave birth to a variety of sabre-like weapons, intended for many tasks. In the following centuries, the ideology of Sarmatism as well as the Polish fascination with Oriental cultures, customs, cuisine and warfare resulted in the szabla becoming an indispensable part of traditional Polish culture.
Modern use
The sabre saw extensive military use in the early 19th century, particularly in the Napoleonic Wars, during which Napoleon used heavy cavalry charges to great effect against his enemies. Shorter versions of the sabre were also used as sidearms by dismounted units, although these were gradually replaced by fascine knives and sword bayonets as the century went on. Although there was extensive debate over the effectiveness of weapons such as the sabre and lance, the sabre remained the standard weapon of cavalry for mounted action in most armies until World War I and in a few armies until World War II. Thereafter it was gradually relegated to the status of a ceremonial weapon, and most horse cavalry was replaced by armoured cavalry from the 1930s onward.
Where horse-mounted cavalry survived into World War II it was generally as mounted infantry without sabres. However the sabre was still carried by German cavalry until after the Polish campaign of 1939, after which this historic weapon was put into storage in 1941.
Romanian cavalry continued to carry their straight "thrusting" sabres on active service until at least 1941.
Napoleonic era
Sabres were commonly used by the British in the Napoleonic era for light cavalry and infantry officers, as well as others. The elegant but effective 1803 pattern sword that the British Government authorized for use by infantry officers during the wars against Napoleon featured a curved sabre blade which was often blued and engraved by the owner in accordance with his personal taste, and was based on the famously agile 1796 light cavalry sabre that was renowned for its brutal cutting power. Sabres were commonly used throughout this era by all armies, in much the same way that the British did.
The popularity of the sabre had rapidly increased in Britain throughout the 18th century for both infantry and cavalry use. This influence was predominately from southern and eastern Europe, with the Hungarians and Austrians listed as sources of influence for the sword and style of swordsmanship in British sources. The popularity of sabres had spread rapidly through Europe in the 16th and 17th centuries, and finally came to dominance as a military weapon in the British army in the 18th century, though straight blades remained in use by some, such as heavy cavalry units. (These were also replaced by sabres soon after the Napoleonic era.)
The introduction of 'pattern' swords in the British army in 1788 led to a brief departure from the sabre in infantry use (though not for light cavalry), in favour of the lighter and straight bladed spadroon. The spadroon was universally unpopular, and many officers began to unofficially purchase and carry sabres once more. In 1799, the army accepted this under regulation for some units, and in 1803, produced a dedicated pattern of sabre for certain infantry officers (flank, rifle and staff officers). The 1803 pattern quickly saw much more widespread use than the regulation intended due to its effectiveness in combat, and fashionable appeal.
Pattern 1796 light cavalry sabre
The most famous British sabre of the Napoleonic era is the 1796 light cavalry model, used by troopers and officers alike (officers versions can vary a little, but are much the same as the pattern troopers sword). It was in part designed by the famous John Le Marchant, who worked to improve on the previous (1788) design based on his experience with the Austrians and Hungarians. Le Marchant also developed the first official British military sword exercise manual based on this experience, and his light cavalry sabre, and style of swordsmanship went on to heavily influence the training of the infantry and the navy.
The 1796 light cavalry sword was known for its brutal cutting power, easily severing limbs, and leading to the (unsubstantiated) myth that the French put in an official complaint to the British about its ferocity. This sword also saw widespread use with mounted artillery units, and the numerous militia units established in Britain to protect against a potential invasion by Napoleon.
Mameluke swords
Though the sabre had already become very popular in Britain, experience in Egypt did lead to a fashion trend for mameluke sword style blades, a type of Middle Eastern scimitar, by some infantry and cavalry officers. These blades differ from the more typical British ones in that they have more extreme curvatures, in that they are usually not fullered, and in that they taper to a finer point. Mameluke swords also gained some popularity in France as well. Arthur Wellesley, 1st Duke of Wellington, himself carried a mameluke-style sword. In 1831, the 'Mameluke' sword became the pattern sword for British generals, as well as officers of the United States Marine Corps; in this last capacity, it is still in such use at the present time.
United States
The American victory over the rebellious forces in the citadel of Tripoli in 1805, during the First Barbary War, led to the presentation of bejewelled examples of these swords to the senior officers of the US Marines. Officers of the US Marine Corps still use a mameluke-pattern dress sword. Although some genuine Turkish kilij sabres were used by Westerners, most "mameluke sabres" were manufactured in Europe; although their hilts were very similar in form to the Ottoman prototype, their blades, even when an expanded was incorporated, tended to be longer, narrower and less curved than those of the true kilij.
In the American Civil War, the sabre was used infrequently as a weapon, but saw notable deployment in the Battle of Brandy Station and at East Cavalry Field at the Battle of Gettysburg in 1863. Many cavalrymen—particularly on the Confederate side—eventually abandoned the long, heavy weapons in favour of revolvers and carbines.
The last sabre issued to US cavalry was the Patton saber of 1913, designed to be mounted to the cavalryman's saddle. The Patton saber is only a saber in name as it is a straight, thrust-centric sword. A US War Department circular dated 18 April 1934 announced that the saber would no longer be issued to cavalry, and that it was to be completely discarded for use as a weapon. Only dress sabers, for use by officers only, and strictly as a badge of rank, were to be retained.
Police
During the 19th and into the early 20th century, sabres were also used by both mounted and dismounted personnel in some European police forces. When the sabre was used by mounted police against crowds, the results could be devastating, as portrayed in a key scene in Doctor Zhivago. The sabre was later phased out in favour of the baton, or nightstick, for both practical and humanitarian reasons. The Gendarmerie of Belgium used them until at least 1950, and the Swedish police forces until 1965.
Dress uniform
Swords with sabre blades remain a component of the dress uniforms worn by most national army, navy, air force, marine and coast guard officers. Some militaries also issue ceremonial swords to their highest-ranking non-commissioned officers; this is seen as an honour since, typically, non-commissioned, enlisted/other-rank military service members are instead issued a cutlass blade rather than a sabre. Swords in the modern military are no longer used as weapons, and serve only ornamental or ceremonial functions. One distinctive modern use of sabres is in the sabre arch, performed for servicemen or women getting married.
Modern sport fencing
The modern fencing sabre bears little resemblance to the cavalry sabre, having a thin, long straight blade. Rather, it is based upon the Italian dueling saber of classical fencing. One of the three weapons used in the sport of fencing, it is a very fast-paced weapon with bouts characterized by quick footwork and cutting with the edge. The valid target area is from the waist up excluding the hands.
The concept of attacking above the waist only is a 20th-century change to the sport; previously sabreurs used to pad their legs against cutting slashes from their opponents. The reason for the above waist rule is unknown, as the sport of sabre fencing is based on the use of infantry sabres, not cavalry sabres.
In recent years, Saber fencing has been developing in Historical European Martial Arts, with blades that closely resemble the historical types, with techniques based on historical records.
| Technology | Swords | null |
75982 | https://en.wikipedia.org/wiki/Defecation | Defecation | Defecation (or defaecation) follows digestion, and is a necessary process by which organisms eliminate a solid, semisolid, or liquid waste material known as feces from the digestive tract via the anus or cloaca. The act has a variety of names ranging from the common, like pooping or crapping, to the technical, e.g. bowel movement, to the obscene (shitting), to the euphemistic ("doing number two", "dropping a deuce" or "taking a dump"), to the juvenile ("making doo-doo"). The topic, usually avoided in polite company, can become the basis for some potty humor.
Humans expel feces with a frequency varying from a few times daily to a few times weekly. Waves of muscular contraction (known as peristalsis) in the walls of the colon move fecal matter through the digestive tract towards the rectum. Undigested food may also be expelled within the feces, in a process called egestion. When birds defecate, they also expel urine and urates in the same mass, whereas other animals may also urinate at the same time, but spatially separated. Defecation may also accompany childbirth and death. Babies defecate a unique substance called meconium prior to eating external foods.
There are a number of medical conditions associated with defecation, such as diarrhea and constipation, some of which can be serious. The feces expelled can carry diseases, most often through the contamination of food. E. coli is a particular concern.
Before toilet training, human feces are most often collected into a diaper. Thereafter, in many societies people commonly defecate into a toilet. However, open defecation, the practice of defecating outside without using a toilet of any kind, is still widespread in some developing countries. Some people defecate into the ocean. First world countries use sewage treatment plants and/or on-site treatment.
Description
Physiology
The rectum ampulla stores fecal waste (also called stool) before it is excreted. As the waste fills the rectum and expands the rectal walls, stretch receptors in the rectal walls stimulate the desire to defecate. This urge to defecate arises from the reflex contraction of rectal muscles, relaxation of the internal anal sphincter, and an initial contraction of the skeletal muscle of the external anal sphincter. If the urge is not acted upon, the material in the rectum is often returned to the colon by reverse peristalsis, where more water is absorbed and the feces are stored until the next mass peristaltic movement of the transverse and descending colon.
When the rectum is full, an increase in pressure within the rectum forces apart the walls of the anal canal, allowing the fecal matter to enter the canal. The rectum shortens as material is forced into the anal canal and peristaltic waves push the feces out of the rectum. The internal and external anal sphincters along with the puborectalis muscle allow the feces to be passed by muscles pulling the anus up over the exiting feces.
Voluntary and involuntary control
The external anal sphincter is under voluntary control whereas the internal anal sphincter is involuntary. In infants, the defecation occurs by reflex action without the voluntary control of the external anal sphincter. Defecation is voluntary in adults. Young children learn voluntary control through the process of toilet training. Once trained, loss of control, called fecal incontinence, may be caused by physical injury, nerve injury, prior surgeries (such as an episiotomy), constipation, diarrhea, loss of storage capacity in the rectum, intense fright, inflammatory bowel disease, psychological or neurological factors, childbirth, or death.
Sometimes, due to the inability to control one's bowel movement or due to excessive fear, defecation (usually accompanied by urination) occurs involuntarily, soiling a person's undergarments. This may cause significant embarrassment to the person if this occurs in the presence of other people or a public place.
Posture
The positions and modalities of defecation are culture-dependent. Squat toilets are used by the vast majority of the world, including most of Africa, Asia, and the Middle East. The use of sit-down toilets in the Western world is a relatively recent development, beginning in the 19th century with the advent of indoor plumbing.
Disease
Regular bowel movements determine the functionality and the health of the alimentary tracts in human body. Defecation is the most common regular bowel movement which eliminates waste from the human body. The frequency of defecation is hard to identify, which can vary from daily to weekly depending on individual bowel habits, the impact from the environment and genetic. If defecation is delayed for a prolonged period the fecal matter may harden, resulting in constipation. If defecation occurs too fast, before excess liquid is absorbed, diarrhea may occur. Other associated symptoms can include abdominal bloating, abdominal pain, and abdominal distention. Disorders of the bowel can seriously impact quality of life and daily activities. The causes of functional bowel disorder are multifactorial, and dietary habits such as food intolerance and low fiber diet are considered to be the primary factors.
Constipation
Constipation, also known as defecatory dysfunction, is difficulty experienced when passing stools. It is one of the most notable alimentary disorders that affects different age groups in the population. Common constipation is associated with abdominal distention, pain or bloating. Research has revealed that chronic constipation complied with higher risk of cardiovascular events such as coronary heart disease and ischemic stroke, while associating with an increasing risk of mortality. Besides dietary factors, psychological traumas and 'pelvic floor disorders' can also cause chronic constipation and defecatory disorder respectively. Multiple interventions, including physical activities, 'high-fibre diet', probiotics and drug therapies can be widely and efficiently used to treat constipation and defecatory disorder.
Inflammatory bowel diseases
Inflammatory disease is characterized as long-lasting, chronic inflammation throughout the gastrointestinal tract. Crohn's disease (CD) and ulcerative colitis (UC) are two universal types of inflammatory bowel disease that have been studied over a century. They are closely related to different environmental risk factors, family genetics, and lifestyle choices such as smoking. Crohn's disease has been found to be related to immune disorders particularly. Different levels of cumulative intestinal injuries can cause different complications, such as fistulae, damage of bowel function, symptom recurrence, disability, etc. Patients can be children or adults. Recent research shows that immunodeficiency and monogenic disorders are the causes in young patients with inflammatory bowel diseases.
Common symptoms of inflammatory bowel diseases differ by the infection level, but may include severe abdominal pain, diarrhea, fatigue, and unexpected weight loss. Crohn's disease can lead to infection of any part of the digestive tract, including ileum to anus. Internal manifestations include diarrhea, abdomen pain, fever, chronic anaemia, etc. External manifestations include impact on skin, joints, eyes, and liver. Significantly reduced microbiota diversity inside the gastrointestinal tract can also be observed. Ulcerative colitis mainly affects the function of the large bowel, and its incidence rate is three times greater than that of Crohn's disease. In terms of clinical features, over 90 percent of patients exhibit constant diarrhea, rectal bleeding, softer stool, mucus in the stool, tenesmus, and abdominal pain. The symptoms may continue for around 6 weeks or even longer.
The inflammatory bowel diseases could be effectively treated by 'pharmacotherapies' to relieve and maintain the symptoms, which showed in 'mucosal healing' and symptoms elimination. However, an optimal therapy for curing both inflammatory diseases are still under research due to the heterogeneity in clinical feature. Although both UC and CD are sharing similar symptoms, the medical treatment of them are distinctively different. Dietary treatment can benefit for curing CD by increase the dietary zinc and fish intake, which is related to mucosal healing of the bowel. Treatments vary from drug treatment to surgery based on the active level of the CD. UC can also be relieved by using immunosuppressive therapy for mild to moderate disease level and application of biological agents for severe cases.
Irritable bowel syndrome
Irritable bowel syndrome is diagnosed as an intestinal disorder with chronic abdominal pain and inconsistency in form of stool, and is a common bowel disease that can be easily diagnosed in modern society. The variation in incident rate can be explained by different diagnostic criteria in different countries, with the 18–34 age group being recognized as the high frequency incident group. The definite cause of irritable bowel syndrome remains a mystery; however, it has been found to relate to multiple factors, such as 'alternation of mood and pressure, sleep disorders, food triggers, changing of dysbiosis and even sexual dysfunction'. One third of irritable bowel syndrome patients has family history with the disease suggesting that genetic predisposition could be a significant cause for irritable bowel syndrome.
Patients with irritable bowel syndrome commonly experience abdominal pain, changes to stool form, recurrent abdominal bloating and gas, co-morbid disorders and alternation in bowel habits that caused diarrhea or constipation. However, anxiety and tension can also be detected, although patients with irritable bowel disease seem healthy. Apart from these typical symptoms, rectal bleeding, unexpected weight loss and increased inflammatory markers require further medical examination and investigation.
Treatment for irritable bowel disease is multimodal. Dietary intervention and pharmacotherapies can both relieve the symptoms to a certain degree. Avoiding allergic food groups can be beneficial by reducing fermentation in the digestive tract and gas production, hence effectively alleviating abdominal pain and bloating. Drug interventions, such as laxatives, loperamide, and lubiprostone are applied to relieve intense symptoms including diarrhea, abdominal pain and constipation. Psychological treatment, dietary supplements and gut-focused hypnotherapy are recommended for targeting depression, mood disorders and sleep disturbance.
Bowel obstruction
Bowel obstruction is a bowel condition which is a blockage that can be found in both the small intestines and large intestines. Increase of contractions can relieve blockages; however, continuous contractions with decreasing functionality may lead to terminated mobility of the small intestines, which then forms the obstruction. At the same time, the lack of contractility encourages liquid and gas accumulation. and "electrolyte disturbances". Small bowel obstruction can result in severe renal damage and hypovolemia. while evolving into "mucosal ischemia and perforation". Patients with small bowel obstruction were found to experience constipation, strangulation and abdominal pain and vomiting. Surgical intervention is primarily used to cure severe small bowel obstruction condition. Nonoperative therapy included nasogastric tube decompression, water-soluble-contrast medium process or symptomatic management can be applied to treat less severe symptoms
According to research, large bowel obstruction is less common than small bowel obstruction, but is still associated with a high mortality rate. Large bowel obstruction, also known as colonic obstruction, includes acute colonic obstruction, where a blockage is formed in the colon. Colonic obstructions frequently occur within the elder population, often accompanied by significant 'comorbidities'. Although colonic malignancy is revealed as the major cause of the colonic obstruction, volvulus has also been founded as a secondary common cause around the world. In addition, lower mobility, unhealthy mentality and restricted living environment are also listed as risk factors. Surgery and colonic stent placements are widely applied for curing colonic obstructions.
Other
Attempting forced expiration of breath against a closed airway (the Valsalva maneuver) is sometimes practiced to induce defecation while on a toilet. This contraction of expiratory chest muscles, diaphragm, abdominal wall muscles, and pelvic diaphragm exerts pressure on the digestive tract. Ventilation at this point temporarily ceases as the lungs push the chest diaphragm down to exert the pressure. Cardiac arrest and other cardiovascular complications can in rare cases occur due to attempting to defecate using the Valsalva maneuver. Valsalva retinopathy is another pathological syndrome associated with the Valsalva maneuver. Thoracic blood pressure rises and as a reflex response the amount of blood pumped by the heart decreases. Death has been known to occur in cases where defecation causes the blood pressure to rise enough to cause the rupture of an aneurysm or to dislodge blood clots (see thrombosis). Also, in releasing the Valsalva maneuver blood pressure falls; this, coupled with standing up quickly to leave the toilet, can result in a blackout.
Society and culture
Open defecation
Open defecation is the human practice of defecating outside (in the open environment) rather than into a toilet. People may choose fields, bushes, forests, ditches, streets, canals or other open space for defecation. They do so because either they do not have a toilet readily accessible or due to traditional cultural practices. The practice is common where sanitation infrastructure and services are not available. Even if toilets are available, behavior change efforts may still be needed to promote the use of toilets.
Open defecation can pollute the environment and cause health problems. High levels of open defecation are linked to high child mortality, poor nutrition, poverty, and large disparities between rich and poor.
Ending open defecation is an indicator being used to measure progress towards the Sustainable Development Goal Number 6. Extreme poverty and lack of sanitation are statistically linked. Therefore, eliminating open defecation is thought to be an important part of the effort to eliminate poverty.
Anal cleansing after defecation
The anus and buttocks may be cleansed after defecation with toilet paper, similar paper products, or other absorbent material. In many cultures, such as Hindu and Muslim, water is used for anal cleansing after defecation, either in addition to using toilet paper or exclusively. When water is used for anal cleansing after defecation, toilet paper may be used for drying the area afterwards. Some doctors and people who work in the science and hygiene fields have stated that switching to using a bidet as a form of anal cleansing after defecation is both more hygienic and more environmentally friendly.
Mythology and tradition
Some peoples have culturally significant stories in which defecation plays a role. For example:
In an Alune and Wemale legend from the island of Seram, Maluku Province, Indonesia, the mythical girl Hainuwele defecates valuable objects.
One of the traditions of Catalonia (Spain) relates to the caganer, a figurine depicting the act of defecation which appears in nativity scenes in Catalonia and neighbouring areas with Catalan culture. The exact origin of the caganer is lost, but the tradition has existed since at least the 18th century.
Psychology
Some aspects of psychology surround the act of defecation. There is an inherent desire for privacy among humans. Freud stipulated a second stage of development, the Anal Stage, which centers around the release of waste from the bladder and bowels. He categorized two types: anal retentive and anal expulsive.
| Biology and health sciences | Basics | Biology |
76046 | https://en.wikipedia.org/wiki/Spanish%20moss | Spanish moss | Spanish moss (Tillandsia usneoides) is an epiphytic flowering plant that often grows upon large trees in tropical and subtropical climates. It is native to much of Mexico, Bermuda, the Bahamas, Central America, South America (as far south as northern Patagonia), the Southern United States, and West Indies. It has been naturalized in Queensland (Australia). It is known as "grandpa's beard" in French Polynesia. It has the widest distribution of any bromeliad.
Most known in the United States, it commonly is found on the southern live oak (Quercus virginiana) and bald cypress (Taxodium distichum) in the lowlands, swamps, and marshes of the mid-Atlantic and Southeastern states, from the coast of southeastern Virginia to Florida and west to southern Arkansas and Texas. While it superficially resembles its namesake, the lichen Usnea, it is neither a lichen nor a moss (instead being a member of the bromeliad family, Bromeliaceae), and it is not native to Spain.
Description
Spanish moss consists of one or more slender stems, bearing alternate thin, curved or curly, and heavily scaled leaves long and broad, that grow vegetatively in a chain-like fashion (pendant), forming hanging structures of up to . The gray-green garlands have occasionally been found hanging down as much as 26 feet (eight meters).
The plant has no roots. Its flowers are yellow-green and small, with spreading petals. The scape is partly hidden within the leaf sheath. Spanish moss propagates both by seed and vegetatively by fragments that are carried on the wind and stick to tree limbs or that are carried to other locations by birds as nesting material.
Taxonomy
Spanish moss is in the family Bromeliaceae (the bromeliads). Formerly, it was placed in the genera Anoplophytum, Caraguata, and Renealmia. The specific name of the plant, usneoides, means "resembling Usnea", a lichen.
Habitat and distribution
Spanish moss' primary range is in the Southeastern United States (including Puerto Rico and the U.S. Virgin Islands) to Argentina, where the climate is warm enough and a relatively high average humidity occurs. In North America, it occurs in a broad band following the Gulf of Mexico and the southern Atlantic coast. The northern limit of its natural range is Northampton County, Virginia, with colonial-era reports of it in southern Maryland, where no populations are now known to exist.
It has been introduced to locations around the world with similar conditions, including Hawaii, where it first established itself in the nineteenth century.
Ecology
Spanish moss is not parasitic: it is an epiphyte that absorbs nutrients and water through its own leaves from the air and rain falling upon it. While its presence rarely kills the trees on which it grows, it occasionally becomes so thick that, by shading the leaves of the tree, it slows the growth rate of the tree. It can use the water-conserving strategy of crassulacean acid metabolism for photosynthesis.
In the southern U.S., the plant seems to show preferences for southern live oak (Quercus virginiana) and bald cypress (Taxodium distichum) because of their high rates of foliar mineral leaching (calcium, magnesium, potassium, and phosphorus) that provides an abundant supply of nutrients to the epiphytic plant. It can also colonize other tree species such as sweetgum (Liquidambar styraciflua), crepe-myrtles (Lagerstroemia spp.), other oaks, and even pines. It also grows more uncommonly on artificial structures such as fencing and telephone lines.
Spanish moss shelters a number of creatures, including rat snakes and three species of bats. One species of jumping spider, Pelegrina tillandsiae, has been found only on Spanish moss. Although widely presumed to infest Spanish moss, in one study of the ecology of the plant, chiggers were not present among thousands of other arthropods identified on the plant.
Spanish moss is sensitive to airborne contaminants. It does not grow in areas where smoke is common, such as near chimneys. It has receded from urban areas due to increasing air pollution.
Culture and folklore
Spanish moss is often associated with Southern Gothic imagery and Deep South culture, due to its propensity for growing in subtropical humid southern locales such as Alabama, Southern Arkansas, Florida, Georgia, Louisiana, Mississippi, North Carolina, South Carolina, east and south Texas, and extreme southern Virginia.
One anecdote about the origin of Spanish moss is called "the Meanest Man Who Ever Lived", in which the man's white hair grew very long and got caught on trees.
Spanish moss was introduced to Hawaii in the nineteenth century. It became a popular ornamental and lei plant. In Hawaii, it was named "ʻumiʻumi-o-Dole" after the beard of Sanford B. Dole, the first president of the Provisional Government of Hawaii. It is also known as hinahina, ("silvery") borrowing the name of the native heliotrope used in lei until shoreline development made access difficult. It has become a substitute for the native hinahina in lei used for pageantry. In the early 21st century the plant was heavily marketed as "Pele's hair"/"lauoho-o-Pele", which actually refers to a type of filamentous volcanic glass.
Human uses
Spanish moss has been used for various purposes, including building insulation, mulch, packing material, mattress stuffing, and fiber. In the early 1900s it was used commercially in the padding of car seats. More than 10,000 tons of processed Spanish moss was produced in 1939. Today, it is collected in smaller quantities for use in arts and crafts, as bedding for flower gardens, and as an ingredient in bousillage, a traditional wall covering material. In some parts of Latin America and Louisiana, it is used in nativity scenes.
In the desert regions of southwestern United States, dried Spanish moss is sometimes used in the manufacture of evaporative coolers, colloquially known as "swamp coolers" (and in some areas as "desert coolers"), which are used to cool homes and offices much less expensively than air conditioners. The cooling technology uses a pump that squirts water onto a pad made of Spanish moss plants; a fan then pulls air through the pad, and into the building. Evaporation of the water on the pads serves to reduce air temperature, cooling the building.
Varieties and cultivars
Tillandsia 'Maurice's Robusta'
Tillandsia 'Munro's Filiformis' – a natural variety with very fine, green leaves that is native to Paraguay and that is also known in the United States by the trade designations Tillandsia usneoides and , it conforms to the description of the now-defunct variety Tillandsia usneoides var. filiformis (André) Mez
Tillandsia 'Odin's Genuina' : a natural variety with brown rather than green or yellow flower petals that is native to Guatemala and Mexico
Tillandsia 'Spanish Gold'
Tillandsia 'Tight and Curly'
Hybrids
Tillandsia 'Nezley' (Tillandsia usneoides × mallemontii)
Tillandsia 'Kimberly' (Tillandsia usneoides × recurvata)
Tillandsia 'Old Man's Gold' (Tillandsia crocata × usneoides)
| Biology and health sciences | Poales | null |
76084 | https://en.wikipedia.org/wiki/Potentiometer | Potentiometer | A potentiometer is a three-terminal resistor with a sliding or rotating contact that forms an adjustable voltage divider. If only two terminals are used, one end and the wiper, it acts as a variable resistor or rheostat.
The measuring instrument called a potentiometer is essentially a voltage divider used for measuring electric potential (voltage); the component is an implementation of the same principle, hence its name.
Potentiometers are commonly used to control electrical devices such as volume controls on audio equipment. It is also used in speed control of fans. Potentiometers operated by a mechanism can be used as position transducers, for example, in a joystick. Potentiometers are rarely used to directly control significant power (more than a watt), since the power dissipated in the potentiometer would be comparable to the power in the controlled load.
Nomenclature
Some terms in the electronics industry used to describe certain types of potentiometers are:
Pot: abbreviation for potentiometer.
Slide pot or slider pot: a potentiometer that is adjusted by sliding the wiper left or right (or up and down, depending on the installation), usually with a finger or thumb
thumb pot or thumbwheel pot: a small rotating potentiometer meant to be adjusted infrequently by means of a small thumbwheel
trimpot or trimmer pot: a trimmer potentiometer typically meant to be adjusted once or infrequently for "fine-tuning" an electrical signal
Construction
Potentiometers consist of a resistive element, a sliding contact (wiper) that moves along the element, making good electrical contact with one part of it, electrical terminals at each end of the element, a mechanism that moves the wiper from one end to the other, and a housing containing the element and wiper.
Many inexpensive potentiometers are constructed with a resistive element (B in cutaway drawing) formed into an arc of a circle usually a little less than a full turn and a wiper (C) sliding on this element when rotated, making electrical contact. The resistive element can be flat or angled. Each end of the resistive element is connected to a terminal (E, G) on the case. The wiper is connected to a third terminal (F), usually between the other two. On panel potentiometers, the wiper is usually the center terminal of three. For single-turn potentiometers, this wiper typically travels just under one revolution around the contact. The only point of ingress for contamination is the narrow space between the shaft and the housing it rotates in.
Another type is the linear slider potentiometer, which has a wiper which slides along a linear element instead of rotating. Contamination can potentially enter anywhere along the slot the slider moves in, making effective sealing more difficult and compromising long-term reliability. An advantage of the slider potentiometer is that the slider position gives a visual indication of its setting. While the setting of a rotary potentiometer can be seen by the position of a marking on the knob, an array of sliders can give a visual impression of settings as in a graphic equalizer or faders on a mixing console.
The resistive element of inexpensive potentiometers is often made of graphite. Other materials used include resistance wire, carbon particles in plastic, and a ceramic/metal mixture called cermet.
Conductive track potentiometers use conductive polymer resistor pastes that contain hard-wearing resins and polymers, solvents, and lubricant, in addition to the carbon that provides the conductive properties.
Multiturn potentiometers are also operated by rotating a shaft, but by several turns rather than less than a full turn. Some multiturn potentiometers have a linear resistive element with a sliding contact moved by a lead screw; others have a helical resistive element and a wiper that turns through 10, 20, or more complete revolutions, moving along the helix as it rotates. Multiturn potentiometers, both user-accessible and preset, allow finer adjustments; rotation through the same angle changes the setting by typically a tenth as much as for a simple rotary potentiometer.
A string potentiometer is a multi-turn potentiometer operated by an attached reel of wire turning against a spring, allowing it to convert linear position to a variable resistance.
User-accessible rotary potentiometers can be fitted with a switch which operates usually at the anti-clockwise extreme of rotation. Before digital electronics became the norm such a component was used to allow radio and television receivers and other equipment to be switched on at minimum volume with an audible click, then the volume increased by turning the same knob. Multiple resistance elements can be ganged together with their sliding contacts on the same shaft, for example in stereo audio amplifiers for volume control. In other applications, such as domestic light dimmers, the normal usage pattern is best satisfied if the potentiometer remains set at its current position, so the switch is operated by a push action, alternately on and off, by axial presses of the knob.
Other potentiometers are enclosed within the equipment and are intended to only be adjusted when calibrating the equipment during manufacture or repair, and not otherwise touched. They are usually physically much smaller than user-accessible potentiometers, and may need to be operated by a screwdriver rather than having a knob. They are usually called "trimmer", "trim[ming]", or "preset" potentiometers (or pots), or the genericized brand name "trimpot".
Resistance–position relationship: "taper"
The relationship between slider position and resistance, known as the "taper" or "law", can be controlled during manufacture by changing the composition or thickness of the resistance coating along the resistance element. Although in principle any taper is possible, two types are widely manufactured: linear and logarithmic (aka "audio taper") potentiometers.
A letter code may be used to identify which taper is used, but the letter code definitions are not standardized. Potentiometers made in Asia and the US are usually marked with an "A" for logarithmic taper or a "B" for linear taper; "C" for the rarely seen reverse logarithmic taper. Others, particularly those from Europe, may be marked with an "A" for linear taper, a "C" or "B" for logarithmic taper, or an "F" for reverse logarithmic taper. The code used also varies between different manufacturers. When a percentage is referenced with a non-linear taper, it relates to the resistance value at the midpoint of the shaft rotation. A 10% log taper would therefore measure 10% of the total resistance at the midpoint of the rotation; i.e. 10% log taper on a 10 kOhm potentiometer would yield 1 kOhm at the midpoint. The higher the percentage, the steeper the log curve.
Linear taper potentiometer
A linear taper potentiometer (linear describes the electrical characteristic of the device, not the geometry of the resistive element) has a resistive element of constant cross-section, resulting in a device where the resistance between the contact (wiper) and one end terminal is proportional to the distance between them. Linear taper potentiometers are used when the division ratio of the potentiometer must be proportional to the angle of shaft rotation (or slider position), for example, controls used for adjusting the centering of the display on an analog cathode-ray oscilloscope. Precision potentiometers have an accurate relationship between resistance and slider position.
Logarithmic potentiometer
A logarithmic taper potentiometer is a potentiometer that has a bias built into the resistive element. Basically this means the center position of the potentiometer is not one half of the total value of the potentiometer. The resistive element is designed to follow a logarithmic taper, aka a mathematical exponent or "squared" profile.
A logarithmic taper potentiometer is constructed with a resistive element that either "tapers" in from one end to the other, or is made from a material whose resistivity varies from one end to the other. This results in a device where output voltage is a logarithmic function of the slider position.
Most (cheaper) "log" potentiometers are not accurately logarithmic, but use two regions of different resistance (but constant resistivity) to approximate a logarithmic law. The two resistive tracks overlap at approximately 50% of the potentiometer rotation; this gives a stepwise logarithmic taper. A logarithmic potentiometer can also be simulated with a linear one and an external resistor. True logarithmic potentiometers are significantly more expensive.
Logarithmic taper potentiometers are often used for volume or signal level in audio systems, as human perception of audio volume is logarithmic, according to the Weber–Fechner law.
Contactless potentiometer
Unlike mechanical potentiometers, non-contact potentiometers use an optical disk to trigger an infrared sensor, or a magnet to trigger a magnetic sensor (as long as there are other types of sensors, such as capacitive, other types of non-contact potentiometers can probably be built), and then an electronic circuit does the signal processing to provide an output signal that can be analogue or digital.
An example of a non-contact potentiometer can be found with the AS5600 integrated circuit. However, absolute encoders must also use similar principles, although being for industrial use, certainly the cost must be unfeasible for use in domestic appliances.
Rheostat
The most common way to vary the resistance in a circuit continuously is to use a rheostat. Because of the change in resistance, they can also be used to adjust magnitude of current in a circuit. The word rheostat was coined in 1843 by Sir Charles Wheatstone, from the Greek rheos meaning "stream", and - -states (from histanai, "to set, to cause to stand") meaning "setter, regulating device", which is a two-terminal variable resistor. For low-power applications (less than about 1 watt) a three-terminal potentiometer is often used, with one terminal unconnected or connected to the wiper.
Where the rheostat must be rated for higher power (more than about 1 watt), it may be built with a resistance wire wound around a semi-circular insulator, with the wiper sliding from one turn of the wire to the next. Sometimes a rheostat is made from resistance wire wound on a heat-resisting cylinder, with the slider made from a number of metal fingers that grip lightly onto a small portion of the turns of resistance wire. The "fingers" can be moved along the coil of resistance wire by a sliding knob thus changing the "tapping" point. Wire-wound rheostats made with ratings up to several thousand watts are used in applications such as DC motor drives, electric welding controls, or in the controls for generators. The rating of the rheostat is given with the full resistance value and the allowable power dissipation is proportional to the fraction of the total device resistance in circuit. Carbon-pile rheostats are used as load banks for testing automobile batteries and power supplies.
Digital potentiometer
A digital potentiometer (often called digipot) is an electronic component that mimics the functions of analog potentiometers. Through digital input signals, the resistance between two terminals can be adjusted, just as in an analog potentiometer. There are two main functional types: volatile, which lose their set position if power is removed, and are usually designed to initialise at the minimum position, and non-volatile, which retain their set position using a storage mechanism similar to flash memory or EEPROM.
Usage of a digipot is far more complex than that of a simple mechanical potentiometer, and there are many limitations to observe; nevertheless they are widely used, often for factory adjustment and calibration of equipment, especially where the limitations of mechanical potentiometers are problematic. A digipot is generally immune to the effects of moderate long-term mechanical vibration or environmental contamination, to the same extent as other semiconductor devices, and can be secured electronically against unauthorised tampering by protecting the access to its programming inputs by various means.
In equipment which has a microprocessor, FPGA or other functional logic which can store settings and reload them to the "potentiometer" every time the equipment is powered up, a multiplying DAC can be used in place of a digipot, and this can offer higher setting resolution, less drift with temperature, and more operational flexibility.
Membrane potentiometers
A membrane potentiometer uses a conductive membrane that is deformed by a sliding element to contact a resistor voltage divider. Linearity can range from 0.50% to 5% depending on the material, design and manufacturing process. The repeat accuracy is typically between 0.1 mm and 1.0 mm with a theoretically infinite resolution. The service life of these types of potentiometers is typically 1 million to 20 million cycles depending on the materials used during manufacturing and the actuation method; contact and contactless (magnetic) methods are available (to sense position). Many different material variations are available such as PET, FR4, and Kapton. Membrane potentiometer manufacturers offer linear, rotary, and application-specific variations. The linear versions can range from 9 mm to 1000 mm in length and the rotary versions range from 20 to 450 mm in diameter, with each having a height of 0.5 mm. Membrane potentiometers can be used for position sensing.
For touch-screen devices using resistive technology, a two-dimensional membrane potentiometer provides x and y coordinates. The top layer is thin glass spaced close to a neighboring inner layer. The underside of the top layer has a transparent conductive coating; the surface of the layer beneath it has a transparent resistive coating. A finger or stylus deforms the glass to contact the underlying layer. Edges of the resistive layer have conductive contacts.
Locating the contact point is done by applying a voltage to opposite edges, leaving the other two edges temporarily unconnected. The voltage of the top layer provides one coordinate. Disconnecting those two edges, and applying voltage to the other two, formerly unconnected, provides the other coordinate. Alternating rapidly between pairs of edges provides frequent position updates. An analog-to-digital converter provides output data.
Advantages of such sensors are that only five connections to the sensor are needed, and the associated electronics is comparatively simple. Another is that any material that depresses the top layer over a small area works well. A disadvantage is that sufficient force must be applied to make contact. Another is that the sensor requires occasional calibration to match touch location to the underlying display. (Capacitive sensors require no calibration or contact force, only proximity of a finger or other conductive object. However, they are significantly more complex.)
Applications
Potentiometers are rarely used to directly control significant amounts of power (more than a watt or so). Instead they are used to adjust the level of analog signals (for example volume controls audio equipment), and as control inputs for electronic circuits. For example, a light dimmer uses a potentiometer to control the switching of a TRIAC and so indirectly to control the brightness of lamps.
Preset potentiometers are widely used throughout electronics wherever adjustments must be made during manufacturing or servicing.
User-actuated potentiometers are widely used as user controls, and may control a very wide variety of equipment functions. The widespread use of potentiometers in consumer electronics declined in the 1990s, with rotary incremental encoders, up/down push-buttons, and other digital controls now more common. However they remain in many applications, such as volume controls and as position sensors.
Audio control
Low-power potentiometers, both slide and rotary, are used to control audio equipment, changing loudness, frequency attenuation, and other characteristics of audio signals.
The 'log pot', that is, a potentiometer has a resistance, taper, or, "curve" (or law) of a logarithmic (log) form, is used as the volume control in audio power amplifiers, where it is also called an "audio taper pot", because the amplitude response of the human ear is approximately logarithmic. It ensures that on a volume control marked 0 to 10, for example, a setting of 5 sounds subjectively half as loud as a setting of 10. There is also an anti-log pot or reverse audio taper which is simply the reverse of a logarithmic potentiometer. It is almost always used in a ganged configuration with a logarithmic potentiometer, for instance, in an audio balance control.
Potentiometers used in combination with filter networks act as tone controls or equalizers.
In audio systems, the word linear, is sometimes applied in a confusing way to describe slide potentiometers because of the straight line nature of the physical sliding motion. The word linear when applied to a potentiometer regardless of being a slide or rotary type, describes a linear relationship of the pot's position versus the measured value of the pot's tap (wiper or electrical output) pin.
Television
Potentiometers were formerly used to control picture brightness, contrast, and color response. A potentiometer was often used to adjust "vertical hold", which affected the synchronization between the receiver's internal sweep circuit (sometimes a multivibrator) and the received picture signal, along with other things such as audio-video carrier offset, tuning frequency (for push-button sets) and so on. It also helps in frequency modulation of waves.
Motion control
Potentiometers can be used as position feedback devices in order to create closed-loop control, such as in a servomechanism. This method of motion control is the simplest method of measuring the angle or displacement.
Transducers
Potentiometers are also very widely used as a part of displacement transducers because of the simplicity of construction and because they can give a large output signal.
Computation
In analog computers, high precision potentiometers are used to scale intermediate results by desired constant factors, or to set initial conditions for a calculation. A motor-driven potentiometer may be used as a function generator, using a non-linear resistance card to supply approximations to trigonometric functions. For example, the shaft rotation might represent an angle, and the voltage division ratio can be made proportional to the cosine of the angle.
Theory of operation
The potentiometer can be used as a voltage divider to obtain a manually adjustable output voltage at the slider (wiper) from a fixed input voltage applied across the two ends of the potentiometer. This is their most common use.
The voltage across can be calculated by:
If is large compared to the other resistances (like the input to an operational amplifier), the output voltage can be approximated by the simpler equation:
(dividing throughout by and cancelling terms with as denominator)
As an example, assume , , , and
Since the load resistance is large compared to the other resistances, the output voltage will be approximately:
Because of the load resistance, however, it will actually be slightly lower: .
One of the advantages of the potential divider compared to a variable resistor in series with the source is that, while variable resistors have a maximum resistance where some current will always flow, dividers are able to vary the output voltage from maximum () to ground (zero volts) as the wiper moves from one end of the potentiometer to the other. There is, however, always a small amount of contact resistance.
In addition, the load resistance is often not known and therefore simply placing a variable resistor in series with the load could have a negligible effect or an excessive effect, depending on the load.
Failure
Ageing may cause intermittent contact between the resistive track and the wiper as it is rotated. In volume control use this causes crackling.
| Technology | Components | null |
76086 | https://en.wikipedia.org/wiki/Electric%20motor | Electric motor | An electric motor is a machine that converts electrical energy into mechanical energy. Most electric motors operate through the interaction between the motor's magnetic field and electric current in a wire winding to generate force in the form of torque applied on the motor's shaft. An electric generator is mechanically identical to an electric motor, but operates in reverse, converting mechanical energy into electrical energy.
Electric motors can be powered by direct current (DC) sources, such as from batteries or rectifiers, or by alternating current (AC) sources, such as a power grid, inverters or electrical generators.
Electric motors may be classified by considerations such as power source type, construction, application and type of motion output. They can be brushed or brushless, single-phase, two-phase, or three-phase, axial or radial flux, and may be air-cooled or liquid-cooled.
Standardized motors provide power for industrial use. The largest are used for ship propulsion, pipeline compression and pumped-storage applications, with output exceeding 100 megawatts.
Applications include industrial fans, blowers and pumps, machine tools, household appliances, power tools, vehicles, and disk drives. Small motors may be found in electric watches. In certain applications, such as in regenerative braking with traction motors, electric motors can be used in reverse as generators to recover energy that might otherwise be lost as heat and friction.
Electric motors produce linear or rotary force (torque) intended to propel some external mechanism. This makes them a type of actuator. They are generally designed for continuous rotation, or for linear movement over a significant distance compared to its size. Solenoids also convert electrical power to mechanical motion, but over only a limited distance.
Components
An electric motor has two mechanical parts: the rotor, which moves, and the stator, which does not. Electrically, the motor consists of two parts, the field magnets and the armature, one of which is attached to the rotor and the other to the stator. Together they form a magnetic circuit. The magnets create a magnetic field that passes through the armature. These can be electromagnets or permanent magnets. The field magnet is usually on the stator and the armature on the rotor, but these may be reversed.
Rotor
The rotor is the moving part that delivers the mechanical power. The rotor typically holds conductors that carry currents, on which the magnetic field of the stator exerts force to turn the shaft.
Stator
The stator surrounds the rotor, and usually holds field magnets, which are either electromagnets (wire windings around a ferromagnetic iron core) or permanent magnets. These create a magnetic field that passes through the rotor armature, exerting force on the rotor windings. The stator core is made up of many thin metal sheets that are insulated from each other, called laminations. These laminations are made of electrical steel, which has a specified magnetic permeability, hysteresis, and saturation. Laminations reduce losses that would result from induced circulating eddy currents that would flow if a solid core were used. Mains powered AC motors typically immobilize the wires within the windings by impregnating them with varnish in a vacuum. This prevents the wires in the winding from vibrating against each other which would abrade the wire insulation and cause premature failures. Resin-packed motors, used in deep well submersible pumps, washing machines, and air conditioners, encapsulate the stator in plastic resin to prevent corrosion and/or reduce conducted noise.
Gap
An air gap between the stator and rotor allows it to turn. The width of the gap has a significant effect on the motor's electrical characteristics. It is generally made as small as possible, as a large gap weakens performance. Conversely, gaps that are too small may create friction in addition to noise.
Armature
The armature consists of wire windings on a ferromagnetic core. Electric current passing through the wire causes the magnetic field to exert a force (Lorentz force) on it, turning the rotor. Windings are coiled wires, wrapped around a laminated, soft, iron, ferromagnetic core so as to form magnetic poles when energized with current.
Electric machines come in salient- and nonsalient-pole configurations. In a salient-pole motor the rotor and stator ferromagnetic cores have projections called poles that face each other. Wire is wound around each pole below the pole face, which become north or south poles when current flows through the wire. In a nonsalient-pole (distributed field or round-rotor) motor, the ferromagnetic core is a smooth cylinder, with the windings distributed evenly in slots around the circumference. Supplying alternating current in the windings creates poles in the core that rotate continuously. A shaded-pole motor has a winding around part of the pole that delays the phase of the magnetic field for that pole.
Commutator
A commutator is a rotary electrical switch that supplies current to the rotor. It periodically reverses the flow of current in the rotor windings as the shaft rotates. It consists of a cylinder composed of multiple metal contact segments on the armature. Two or more electrical contacts called brushes made of a soft conductive material like carbon press against the commutator. The brushes make sliding contact with successive commutator segments as the rotator turns, supplying current to the rotor. The windings on the rotor are connected to the commutator segments. The commutator reverses the current direction in the rotor windings with each half turn (180°), so the torque applied to the rotor is always in the same direction. Without this reversal, the direction of torque on each rotor winding would reverse with each half turn, stopping the rotor. Commutated motors have been mostly replaced by brushless motors, permanent magnet motors, and induction motors.
Shaft
The motor shaft extends outside of the motor, where it satisfies the load. Because the forces of the load are exerted beyond the outermost bearing, the load is said to be overhung.
Bearings
The rotor is supported by bearings, which allow the rotor to turn on its axis by transferring the force of axial and radial loads from the shaft to the motor housing.
History
Early motors
Before modern electromagnetic motors, experimental motors that worked by electrostatic force were investigated. The first electric motors were simple electrostatic devices described in experiments by Scottish monk Andrew Gordon and American experimenter Benjamin Franklin in the 1740s. The theoretical principle behind them, Coulomb's law, was discovered but not published, by Henry Cavendish in 1771. This law was discovered independently by Charles-Augustin de Coulomb in 1785, who published it so that it is now known by his name. Due to the difficulty of generating the high voltages they required, electrostatic motors were never used for practical purposes.
The invention of the electrochemical battery by Alessandro Volta in 1799 made the production of persistent electric currents possible. Hans Christian Ørsted discovered in 1820 that an electric current creates a magnetic field, which can exert a force on a magnet. It only took a few weeks for André-Marie Ampère to develop the first formulation of the electromagnetic interaction and present the Ampère's force law, that described the production of mechanical force by the interaction of an electric current and a magnetic field.
Michael Faraday gave the first demonstration of the effect with a rotary motion on 3 September 1821 in the basement of the Royal Institution. A free-hanging wire was dipped into a pool of mercury, on which a permanent magnet (PM) was placed. When a current was passed through the wire, the wire rotated around the magnet, showing that the current gave rise to a close circular magnetic field around the wire. Faraday published the results of his discovery in the Quarterly Journal of Science, and sent copies of his paper along with pocket-sized models of his device to colleagues around the world so they could also witness the phenomenon of electromagnetic rotations. This motor is often demonstrated in physics experiments, substituting brine for (toxic) mercury. Barlow's wheel was an early refinement to this Faraday demonstration, although these and similar homopolar motors remained unsuited to practical application until late in the century.
In 1827, Hungarian physicist Ányos Jedlik started experimenting with electromagnetic coils. After Jedlik solved the technical problems of continuous rotation with the invention of the commutator, he called his early devices "electromagnetic self-rotors". Although they were used only for teaching, in 1828 Jedlik demonstrated the first device to contain the three main components of practical DC motors: the stator, rotor and commutator. The device employed no permanent magnets, as the magnetic fields of both the stationary and revolving components were produced solely by the currents flowing through their windings.
DC motors
The first commutator capable of turning machinery was invented by English scientist William Sturgeon in 1832. Following Sturgeon's work, a commutator-type direct-current electric motor was built by American inventors Thomas Davenport and Emily Davenport, which he patented in 1837. The motors ran at up to 600 revolutions per minute, and powered machine tools and a printing press. Due to the high cost of primary battery power, the motors were commercially unsuccessful and bankrupted the Davenports. Several inventors followed Sturgeon in the development of DC motors, but all encountered the same battery cost issues. As no electricity distribution system was available at the time, no practical commercial market emerged for these motors.
After many other more or less successful attempts with relatively weak rotating and reciprocating apparatus Prussian/Russian Moritz von Jacobi created the first real rotating electric motor in May 1834. It developed remarkable mechanical output power. His motor set a world record, which Jacobi improved four years later in September 1838. His second motor was powerful enough to drive a boat with 14 people across a wide river. It was also in 1839/40 that other developers managed to build motors with similar and then higher performance.
In 1827–1828, Jedlik built a device using similar principles to those used in his electromagnetic self-rotors that was capable of useful work. He built a model electric vehicle that same year.
A major turning point came in 1864, when Antonio Pacinotti first described the ring armature (although initially conceived in a DC generator, i.e. a dynamo). This featured symmetrically grouped coils closed upon themselves and connected to the bars of a commutator, the brushes of which delivered practically non-fluctuating current. The first commercially successful DC motors followed the developments by Zénobe Gramme who, in 1871, reinvented Pacinotti's design and adopted some solutions by Werner Siemens.
A benefit to DC machines came from the discovery of the reversibility of the electric machine, which was announced by Siemens in 1867 and observed by Pacinotti in 1869. Gramme accidentally demonstrated it on the occasion of the 1873 Vienna World's Fair, when he connected two such DC devices up to 2 km from each other, using one of them as a generator and the other as motor.
The drum rotor was introduced by Friedrich von Hefner-Alteneck of Siemens & Halske to replace Pacinotti's ring armature in 1872, thus improving the machine efficiency. The laminated rotor was introduced by Siemens & Halske the following year, achieving reduced iron losses and increased induced voltages. In 1880, Jonas Wenström provided the rotor with slots for housing the winding, further increasing the efficiency.
In 1886, Frank Julian Sprague invented the first practical DC motor, a non-sparking device that maintained relatively constant speed under variable loads. Other Sprague electric inventions about this time greatly improved grid electric distribution (prior work done while employed by Thomas Edison), allowed power from electric motors to be returned to the electric grid, provided for electric distribution to trolleys via overhead wires and the trolley pole, and provided control systems for electric operations. This allowed Sprague to use electric motors to invent the first electric trolley system in 1887–88 in Richmond, Virginia, the electric elevator and control system in 1892, and the electric subway with independently powered centrally-controlled cars. The latter were first installed in 1892 in Chicago by the South Side Elevated Railroad, where it became popularly known as the "L". Sprague's motor and related inventions led to an explosion of interest and use in electric motors for industry. The development of electric motors of acceptable efficiency was delayed for several decades by failure to recognize the extreme importance of an air gap between the rotor and stator. Efficient designs have a comparatively small air gap. The St. Louis motor, long used in classrooms to illustrate motor principles, is inefficient for the same reason, as well as appearing nothing like a modern motor.
Electric motors revolutionized industry. Industrial processes were no longer limited by power transmission using line shafts, belts, compressed air or hydraulic pressure. Instead, every machine could be equipped with its own power source, providing easy control at the point of use, and improving power transmission efficiency. Electric motors applied in agriculture eliminated human and animal muscle power from such tasks as handling grain or pumping water. Household uses (sauch as washing machines, dishwashers, fans, air conditioners and refrigerators) of electric motors reduced heavy labor in the home and made higher standards of convenience, comfort and safety possible. Today, electric motors consume more than half of the electric energy produced in the US.
AC motors
In 1824, French physicist François Arago formulated the existence of rotating magnetic fields, termed Arago's rotations, which, by manually turning switches on and off, Walter Baily demonstrated in 1879 as in effect the first primitive induction motor. In the 1880s many inventors were trying to develop workable AC motors because AC's advantages in long-distance high-voltage transmission were offset by the inability to operate motors on AC.
The first alternating-current commutatorless induction motor was invented by Galileo Ferraris in 1885. Ferraris was able to improve his first design by producing more advanced setups in 1886. In 1888, the Royal Academy of Science of Turin published Ferraris's research detailing the foundations of motor operation, while concluding at that time that "the apparatus based on that principle could not be of any commercial importance as motor."
Possible industrial development was envisioned by Nikola Tesla, who invented independently his induction motor in 1887 and obtained a patent in May 1888. In the same year, Tesla presented his paper A New System of Alternate Current Motors and Transformers to the AIEE that described three patented two-phase four-stator-pole motor types: one with a four-pole rotor forming a non-self-starting reluctance motor, another with a wound rotor forming a self-starting induction motor, and the third a true synchronous motor with separately excited DC supply to rotor winding. One of the patents Tesla filed in 1887, however, also described a shorted-winding-rotor induction motor. George Westinghouse, who had already acquired rights from Ferraris (US$1,000), promptly bought Tesla's patents (US$60,000 plus US$2.50 per sold hp, paid until 1897), employed Tesla to develop his motors, and assigned C.F. Scott to help Tesla; however, Tesla left for other pursuits in 1889. The constant speed AC induction motor was found not to be suitable for street cars, but Westinghouse engineers successfully adapted it to power a mining operation in Telluride, Colorado in 1891. Westinghouse achieved its first practical induction motor in 1892 and developed a line of polyphase 60 hertz induction motors in 1893, but these early Westinghouse motors were two-phase motors with wound rotors. B.G. Lamme later developed a rotating bar winding rotor.
Steadfast in his promotion of three-phase development, Mikhail Dolivo-Dobrovolsky invented the three-phase induction motor in 1889, of both types cage-rotor and wound rotor with a starting rheostat, and the three-limb transformer in 1890. After an agreement between AEG and Maschinenfabrik Oerlikon, Doliwo-Dobrowolski and Charles Eugene Lancelot Brown developed larger models, namely a 20-hp squirrel cage and a 100-hp wound rotor with a starting rheostat. These were the first three-phase asynchronous motors suitable for practical operation. Since 1889, similar developments of three-phase machinery were started Wenström. At the 1891 Frankfurt International Electrotechnical Exhibition, the first long distance three-phase system was successfully presented. It was rated 15 kV and extended over 175 km from the Lauffen waterfall on the Neckar river. The Lauffen power station included a 240 kW 86 V 40 Hz alternator and a step-up transformer while at the exhibition a step-down transformer fed a 100-hp three-phase induction motor that powered an artificial waterfall, representing the transfer of the original power source. The three-phase induction is now used for the vast majority of commercial applications. Mikhail Dolivo-Dobrovolsky claimed that Tesla's motor was not practical because of two-phase pulsations, which prompted him to persist in his three-phase work.
The General Electric Company began developing three-phase induction motors in 1891. By 1896, General Electric and Westinghouse signed a cross-licensing agreement for the bar-winding-rotor design, later called the squirrel-cage rotor. Induction motor improvements flowing from these inventions and innovations were such that a 100-horsepower induction motor currently has the same mounting dimensions as a 7.5-horsepower motor in 1897.
Twenty-first century
In 2022, electric motor sales were estimated to be 800 million units, increasing by 10% annually. Electric motors consume ≈50% of the world's electricity. Since the 1980s, the market share of DC motors has declined in favor of AC motors.
Inputs
Power supply
A DC motor is usually supplied through a split ring commutator as described above.
AC motors' commutation can be achieved using either a slip ring commutator or external commutation. It can be fixed-speed or variable-speed control type, and can be synchronous or asynchronous. Universal motors can run on either AC or DC.
Control
DC motors can be operated at variable speeds by adjusting the voltage applied to the terminals or by using pulse-width modulation (PWM).
AC motors operated at a fixed speed are generally powered directly from the grid or through motor soft starters.
AC motors operated at variable speeds are powered with various power inverter, variable-frequency drive or electronic commutator technologies.
The term electronic commutator is usually associated with self-commutated brushless DC motor and switched reluctance motor applications.
Types
Electric motors operate on one of three physical principles: magnetism, electrostatics and piezoelectricity.
In magnetic motors, magnetic fields are formed in both the rotor and the stator. The product between these two fields gives rise to a force and thus a torque on the motor shaft. One or both of these fields changes as the rotor turns. This is done by switching the poles on and off at the right time, or varying the strength of the pole.
Motors can be designed to operate on DC current, on AC current, or some types can work on either.
AC motors can be either asynchronous or synchronous. Synchronous motors require the rotor to turn at the same speed as the stator's rotating field. Asynchronous rotors relax this constraint.
A fractional-horsepower motor either has a rating below about 1 horsepower (0.746 kW), or is manufactured with a frame size smaller than a standard 1 HP motor. Many household and industrial motors are in the fractional-horsepower class.
| Technology | Tools and machinery | null |
76097 | https://en.wikipedia.org/wiki/Honeysuckle | Honeysuckle | Honeysuckles are arching shrubs or twining vines in the genus Lonicera () of the family Caprifoliaceae. The genus includes 158 species native to northern latitudes in North America, Eurasia, and North Africa. Widely known species include Lonicera periclymenum (common honeysuckle or woodbine), Lonicera japonica (Japanese honeysuckle, white honeysuckle, or Chinese honeysuckle) and Lonicera sempervirens (coral honeysuckle, trumpet honeysuckle, or woodbine honeysuckle). L. japonica is a highly invasive species considered a significant pest in parts of North America, Europe, South America, New Zealand, Australia, and Africa.
Some species are highly fragrant and colorful, so are cultivated as ornamental garden plants. In North America, hummingbirds are attracted to the flowers, especially L. sempervirens and L. ciliosa (orange honeysuckle). Honeysuckle derives its name from the edible sweet nectar obtainable from its tubular flowers. The name Lonicera stems from Adam Lonicer, a Renaissance botanist.
Description
Most species of Lonicera are hardy twining climbers, with a minority of shrubby habit. Some species (including Lonicera hildebrandiana from the Himalayan foothills and L. etrusca from the Mediterranean) are tender and can be grown outside only in subtropical zones. The leaves are opposite, simple oval, long; most are deciduous but some are evergreen.
Many of the species have sweetly scented, bilaterally symmetrical flowers that produce a sweet, edible nectar, and most flowers are borne in clusters of two (leading to the common name of "twinberry" for certain North American species). Both shrubby and vining sorts have strongly fibrous stems which have been used for binding and textiles.
The fruit is a red, blue or black spherical or elongated berry containing several seeds; in most species the berries are mildly poisonous, but in a few (notably Lonicera caerulea) they are edible and grown for home use and commerce. Most honeysuckle berries are attractive to wildlife, which has led to species such as L. japonica and L. maackii spreading invasively outside of their home ranges. Many species of Lonicera are eaten by the larvae of some Lepidoptera species—see a list of Lepidoptera that feed on honeysuckles.
Invasive species
The spread of L. japonica in North America began in the United States in 1806, and it was widely cultivated by the 1860s. It was first discovered in Canada in Ontario forests in 1976, and became invasive by 2007. L. japonica was introduced in Australia between 1820 and 1840.
Several species of honeysuckle have become invasive when introduced outside their native range, particularly in North America, Europe, South America, Australia, and Africa. Invasive species include L. japonica, L. maackii, L. morrowii, L. tatarica, and the hybrid between the last two, L. × bella.
Cultivation
Honeysuckles are valued as garden plants, for their ability to cover unsightly walls and outbuildings, their profuse tubular flowers in early summer, and the intense fragrance of many varieties. The hardy climbing types need their roots in shade, and their flowering tops in sunlight or very light shade. Varieties need to be chosen with care, as they can become substantial. Cultivars of the dense, small-leaved L. nitida are used as low, narrow hedges.
The following hybrids have gained the Royal Horticultural Society's Award of Garden Merit:
L. × heckrottii 'Gold Flame'
L. 'Mandarin'
L. × purpusii 'Winter Beauty'
L. × tellmanniana
Other cultivars are dealt with under their species names.
The honeysuckle species L. japonica is grown as a commercial crop for traditional Chinese medicine use.
Honeysuckle is also used to scent Chinese teas in a process similar to Jasmine tea. This was popularized in the Qing dynasty.
Phytochemicals
Component analyses of berries from 27 different cultivars and 3 genotypes of edible honeysuckle (Lonicera caerulea var. kamtschatica) showed the presence of iridoids, anthocyanins, flavonols, flavanonols, flavones, flavan-3-ols, and phenolic acids. While sugars determine the level of sweetness in the berries, organic acids and polyphenols are responsible for the sour taste and tartness. Some 51 of the same compounds in berries are found in flowers, although the proportions of these compounds varied among cultivars studied.
Interaction with other species
Many insects in the order Lepidoptera visit honeysuckles as a food source. An example of this is the moth Deilephila elpenor. This nocturnal species of moth is especially attracted to honeysuckles, and they visit the flowers at night to feed on their nectar.
Species
158 species are accepted.
Lonicera acuminata or Lonicera pampaninii – fragrant grove honeysuckle or vine honeysuckle
Lonicera affinis
Lonicera alberti
Lonicera albiflora – white honeysuckle
Lonicera alpigena – alpine honeysuckle
Lonicera altmannii
Lonicera × americana
Lonicera angustifolia
Lonicera annamensis
Lonicera arborea
Lonicera arizonica – Arizona honeysuckle
Lonicera asperifolia
Lonicera × bella – Bell's honeysuckle or showy fly honeysuckle
Lonicera biflora
Lonicera bournei
Lonicera bracteolaris
Lonicera buschiorum
Lonicera caerulea – blue-berried honeysuckle
Lonicera calcarata
Lonicera cambodiana
Lonicera canadensis – Canada fly honeysuckle, American fly honeysuckle
Lonicera caprifolium – goat-leaf honeysuckle, perfoliate honeysuckle
Lonicera caucasica
Lonicera cerasina
Lonicera cerviculata
Lonicera chamissoi
Lonicera chrysantha – Chrysantha honeysuckle
Lonicera ciliosa – orange honeysuckle
Lonicera confusa
Lonicera conjugialis – purpleflower honeysuckle
Lonicera crassifolia
Lonicera cyanocarpa
Lonicera deleiensis
Lonicera demissa
Lonicera dioica – limber honeysuckle
Lonicera elisae
Lonicera etrusca – Etruscan honeysuckle
Lonicera fargesii
Lonicera ferdinandii
Lonicera ferruginea
Lonicera flava – yellow honeysuckle
Lonicera floribunda
Lonicera fragrantissima – winter honeysuckle
Lonicera glabrata
Lonicera glehnii
Lonicera gracilipes
Lonicera griffithii
Lonicera guatemalensis
Lonicera guillonii
Lonicera gynochlamydea
Lonicera harae
Lonicera × heckrottii – golden flame honeysuckle
Lonicera × helvetica
Lonicera heterotricha
Lonicera hildebrandiana – giant Burmese honeysuckle
Lonicera himalayensis
Lonicera hirsuta – hairy honeysuckle
Lonicera hispida
Lonicera hispidula – pink honeysuckle
Lonicera humilis
Lonicera hypoglauca
Lonicera hypoleuca
Lonicera iberica
Lonicera iliensis
Lonicera implexa
Lonicera interrupta – Chaparral honeysuckle
Lonicera involucrata – bearberry honeysuckle
Lonicera × italica
Lonicera japonica – Japanese honeysuckle
Lonicera kansuensis
Lonicera kawakamii
Lonicera korolkowii – blueleaf honeysuckle
Lonicera kurobushiensis
Lonicera laceana
Lonicera lanceolata
Lonicera ligustrina
Lonicera ligustrina var. ligustrina
Lonicera ligustrina var. pileata (syn. Lonicera pileata ) – privet honeysuckle
Lonicera ligustrina var. yunnanensis (syn. Lonicera nitida ) – boxleaf honeysuckle
Lonicera litangensis
Lonicera longiflora
Lonicera longituba
Lonicera maackii – Amur honeysuckle
Lonicera macrantha
Lonicera macranthoides
Lonicera magnibracteata
Lonicera malayana
Lonicera maximowiczii
Lonicera mexicana
Lonicera micrantha
Lonicera microphylla
Lonicera minutifolia
Lonicera mochidzukiana
Lonicera modesta
Lonicera morrowii – Morrow's honeysuckle
Lonicera mucronata
Lonicera myrtilloides
Lonicera nervosa
Lonicera nigra – black-berried honeysuckle
Lonicera nummulariifolia
Lonicera oblata
Lonicera oblongifolia – swamp fly honeysuckle
Lonicera obovata
Lonicera olgae
Lonicera oreodoxa
Lonicera pamirica
Lonicera paradoxa
Lonicera periclymenum – (common) honeysuckle, European honeysuckle, or woodbine
Lonicera pilosa – Mexican honeysuckle
Lonicera praeflorens
Lonicera purpurascens
Lonicera pyrenaica – Pyrenean honeysuckle
Lonicera quinquelocularis – translucent honeysuckle
Lonicera reticulata – grape honeysuckle
Lonicera retusa
Lonicera robertsonii
Lonicera rupicola
Lonicera ruprechtiana – Manchurian honeysuckle
Lonicera × sargentii
Lonicera schmitziana
Lonicera semenovii
Lonicera sempervirens – trumpet honeysuckle
Lonicera setifera
Lonicera siamensis
Lonicera similis – var. delavayi – Delavay honeysuckle
Lonicera sinomacrantha
Lonicera sovetkinae
Lonicera spinosa
Lonicera splendida – evergreen honeysuckle
Lonicera stabiana
Lonicera stephanocarpa
Lonicera steveniana
Lonicera strophiophora
Lonicera subaequalis
Lonicera subhispida
Lonicera sublabiata
Lonicera subsessilis
Lonicera subspicata – southern honeysuckle
Lonicera sumatrana
Lonicera taiwanensis
Lonicera tangutica
Lonicera tatarica – Tatarian honeysuckle
Lonicera tatarinowii
Lonicera tolmatchevii
Lonicera tomentella
Lonicera tragophylla – Chinese honeysuckle
Lonicera tricalysioides
Lonicera trichosantha
Lonicera tschonoskii
Lonicera tubuliflora
Lonicera tulinensis
Lonicera utahensis – Utah honeysuckle
Lonicera uzenensis
Lonicera vaccinioides
Lonicera vidalii
Lonicera villosa – mountain fly honeysuckle
Lonicera webbiana
Lonicera xylosteum – fly woodbine
Lonicera yunnanensis
Lonicera zeravshanica
Several fossil species are known from the Miocene of Asia.
| Biology and health sciences | Dipsacales | Plants |
76098 | https://en.wikipedia.org/wiki/Syringa | Syringa | Syringa is a genus of 12 currently recognized species of flowering woody plants in the olive family or Oleaceae called lilacs. These lilacs are native to woodland and scrub from southeastern Europe to eastern Asia, and widely and commonly cultivated in temperate areas elsewhere.
The genus is most closely related to Ligustrum (privet), classified with it in Oleaceae tribus Oleeae subtribus Ligustrinae.
Lilacs are used as food plants by the larvae of some moth species, including lilac leaf mining moth, privet hawk moth, copper underwing, scalloped oak and Svensson's copper underwing.
Description
They are small trees, ranging in size from tall, with stems up to diameter. The leaves are opposite (occasionally in whorls of three) in arrangement, and their shape is simple and heart-shaped to broad lanceolate in most species, but pinnate in a few species (e.g. S. protolaciniata, S. pinnatifolia).
Flowers
The flowers are produced in spring, each flower being in diameter with a four-lobed corolla, the corolla tube narrow, long; they are monoecious, with fertile stamens and stigma in each flower. The usual flower colour is a shade of purple (often a light purple or "lilac"), but white, pale yellow and pink, and even a dark burgundy color are also found.
The flowers grow in large panicles, and in several species have a strong fragrance. Flowering varies between mid spring to early summer, depending on the species. One particular cultivar, trademark Bloomerang, first blooms in spring and then again late summer through fall.
Fruit
The fruit is a dry, brown capsule, splitting in two at maturity to release the two winged seeds.
Etymology
The English common name "lilac" is from the French lilac via the from meaning the indigo plant or nilak meaning "bluish"; both lilanj and nilak come from Persian nīl "indigo" or nili "dark blue".
Taxonomy
The genus Syringa was first formally described in 1753 by Carl Linnaeus and the description was published in Species Plantarum. The genus name Syringa is derived from Ancient Greek word syrinx meaning "pipe" or "tube" and refers to the hollow branches of S. vulgaris.
Homonym Syringa Tourn. ex Adans. is a heterotypic synonym of Philadelphus.
Cultivation and uses
Lilacs are popular shrubs in parks and gardens throughout the temperate zone, and several hybrids and numerous cultivars have been developed. The term French lilac is often used to refer to modern double-flowered cultivars, thanks to the work of prolific breeder Victor Lemoine. Lilacs grow most successfully in well-drained soils, particularly those based on chalk. They flower on old wood, and produce more flowers if unpruned. If pruned, the plant responds by producing fast-growing young vegetative growth with no flowers, in an attempt to restore the removed branches. Lilac bushes can be prone to powdery mildew disease.
Lilac wood is not commonly used or commercially harvested due to the small size of the tree. It is a relatively hard wood, with an estimated Janka hardness of 2,350 lbf (10,440 N), and is reportedly good for woodturning
The sapwood is typically cream-coloured and the heartwood can have various streaks of brown and purple.
Species have been historically used in various traditional medicines in Asia for treating ailments including cough, diarrhea, acute icteric hepatitis, vomiting, abdominal pain, and bronchitis. Compounds isolated from species of Syringa include phenylpropanoids such as syringin and iridoids such as oleuropein. Substituent compounds, such as iridoids, as well as crude extracts from Syringa plants have been shown to have to have effects including antitumor, antihypertensive, anti-inflammatory, antioxidant, and antifungal activities in pharmacological studies.
Symbolism
Lilacs are often considered to symbolize first love.
In Greece, Macedonia, Lebanon, and Cyprus, the lilac is strongly associated with Easter time because it flowers around that time; it is consequently called paschalia.
In the poem When Lilacs Last in the Dooryard Bloom'd, by Walt Whitman, lilacs are a reference to Abraham Lincoln.
The music-hall song by Ivor Novello, We'll Gather Lilacs, first performed in 1945, speaks of the longing of two lovers to be reunited in a traditional English rural setting. It has since been recorded and performed by numerous artists.
Syringa vulgaris is the state flower of New Hampshire, because it "is symbolic of that hardy character of the men and women of the Granite State."
Festivals
Several locations in North America hold annual Lilac Festivals, including:
The Arnold Arboretum in Boston, Massachusetts, which celebrates "Lilac Sunday" every May. The Arboretum shows off its collection of over 422 lilac plants, of 194 different varieties. Lilac Sunday is the only day of the year when picnicking is allowed on the grounds of the Arboretum.
Lombard, Illinois, called the "Lilac Village", which has an annual lilac festival and parade in May. The village also contains Lilacia Park, a garden with over 200 varieties of lilacs, as well as over 50 kinds of tulips.
Mackinac Island, in Michigan, which celebrates a weeklong lilac festival and lilac parade each June.
Rochester, New York, which has held its Lilac Festival since 1898, hosts the longest-running festival in North America. Held in Highland Park, this celebration features 1,200 shrubs, representing over 500 varieties, many of which were developed in Rochester. It is the largest collection of varieties at any single place.
The Royal Botanical Gardens near Hamilton, Ontario, which holds its Lilac Celebration each May.
Spokane, Washington, known as the "Lilac City", which holds an annual lilac festival and lilac parade.
Franktown, Ontario, Canada, known as the Lilac Capital of Canada, holds an annual festival. With drystone masonry demonstrations and horse pulled wagon rides.
Calgary, Alberta, Canada holds an annual one-day Lilac Festival, which is primarily a street festival.
Species
Species and subspecies currently accepted as of July 2016:
Syringa emodi Wall. ex Royle – Himalayan lilac - northern India, Pakistan, Tibet, Nepal
Syringa josikaea J.Jacq. ex Rchb.f. – Hungarian lilac - Carpathian Mountains of Romania and Ukraine
Syringa komarowii C.K.Schneid. – nodding lilac - Gansu, Hubei, Shaanxi, Sichuan, Yunnan
Syringa oblata Lindl. – early blooming lilac or broadleaf lilac - Korea, Gansu, Hebei, Henan, Jilin, Liaoning, Inner Mongolia, Ningxia, Qinghai, Shaanxi, Shandong, Shanxi, Sichuan
Syringa oblata subsp. dilatata – Korean early lilac - Nakai - Korea, Jilin, LiaoningSyringa pinetorum W.W.Sm. – Sichuan, Tibet, YunnanSyringa pinnatifolia Hemsl. – Gansu, Inner Mongolia, Ningxia, Qinghai, Shaanxi, SichuanSyringa pubescens Turcz. – Korea, Gansu, Hebei, Henan, Hubei, Jilin, Liaoning, Ningxia, Qinghai, Shaanxi, Shandong, Shanxi, SichuanSyringa reticulata (Blume) H.Hara (syn. S. pekinensis) – Japanese tree lilac - Primorye, Japan, Korea, Gansu, Hebei, Heilongjiang, Henan, Jilin, Liaoning, Inner Mongolia, Ningxia, Shaanxi, Shanxi, Sichuan Syringa tomentella Bureau & Franch. – Sichuan, Tibet, YunnanSyringa villosa Vahl – villous lilac - Primorye, Korea, Hebei, Shanxi, Heilongjiang, Jilin, Liaoning Syringa vulgaris L. – common lilac - native to Balkans; naturalized in western and central Europe, and many scattered locations in North America
Hybrids
S. × chinensis (S. vulgaris × S. persica)
S. × diversifolia (S. oblata × S. pinnatifolia)S. × henryi (S. josikaea × S. villosa)S. × hyacinthiflora (S. oblata × S. vulgaris)S. × josiflexa (S. josikaea × S. komarowii)
S. × laciniata (S. protolaciniata × S. vulgaris) – cut-leaf lilac or cutleaf lilac
S. × persica L. (syn Syringa protolaciniata) – Persian lilac - Afghanistan, Pakistan, western Himalayas, Gansu, QinghaiS. × prestoniae (S. komarowii × S. villosa)S. × swegiflexa (S. komarowii × S. sweginzowii'')
Gallery
| Biology and health sciences | Lamiales | Plants |
76099 | https://en.wikipedia.org/wiki/Camellia | Camellia | Camellia (pronounced or ) is a genus of flowering plants in the family Theaceae. They are found in tropical and subtropical areas in eastern and southern Asia, from the Himalayas east to Japan and Indonesia. There are more than 220 described species; almost all are found in southern China and Indochina. Camellias are popular ornamental, tea, and woody-oil plants cultivated worldwide for centuries. Over 26,000 cultivars, with more than 51,000 cultivar names, including synonyms, have been registered or published.
Of economic importance in East Asia, Southeast Asia, and the Indian subcontinent, leaves of C. sinensis are processed to create the popular beverage tea. The ornamental C. japonica, C. sasanqua and their hybrids are the source of hundreds of garden cultivars. C. oleifera produces tea seed oil, used in cooking and cosmetics.
Taxonomy
The genus was named by Linnaeus after the Jesuit botanist Georg Joseph Kamel, who worked in the Philippines and described one of its species (although Linnaeus did not refer to Kamel's account when discussing the genus).
Botany
Camellias are evergreen shrubs or small trees up to tall. Their leaves are alternately arranged, simple, thick, serrated, and usually glossy.
Flowers and fruit
Their flowers are usually large and conspicuous, one to 12 cm in diameter, with five to nine petals in naturally occurring species of camellias. The colors of the flowers vary from white through pink colors to red; truly yellow flowers are found only in South China and Vietnam. Tea varieties are always white-flowered. Camellia flowers throughout the genus are characterized by a dense bouquet of conspicuous yellow stamens, often contrasting with the petal colors. Some research has shown that the colour of petals in some species' flowers indicate their size and how they are pollinated; species with red or yellow flowers are pollinated by sunbirds whereas species with white flowers are smaller in diameter and are pollinated by bees.
The fruit of camellia plants is a dry capsule, sometimes subdivided into up to five compartments. Each compartment contains up to eight seeds.
Growth
The various species of camellia plants are generally well-adapted to acid soils rich in humus, and most species do not grow well on chalky soil or other calcium-rich soils. Most species of camellias also require a large amount of water, either from natural rainfall or from irrigation, and the plants will not tolerate droughts. However, some of the more unusual camellias – typically species from karst soils in Vietnam – can grow without too much water.
Camellia plants usually have a rapid growth rate. Typically, they will grow about 30 cm per year until mature – however, this varies depending on their variety and geographical location.
Ecology
Camellia plants are used as food plants by the larvae of some Lepidoptera species. Leaves of Camellia japonica are susceptible to the fungal parasite Mycelia sterile (see below for the significance), mycelia sterile PF1022 produces a metabolite named PF1022A that is used to produce emodepside, an anthelmintic drug.
Due to habitat destruction, several camellias have become rare in their natural range. One of these is the aforementioned C. reticulata, grown commercially in thousands for horticulture and oil production but rare enough in its natural range to be considered a threatened species.
Use by humans
Camellia sinensis, the tea plant, is of major commercial importance because tea is made from its leaves. The species C. sinensis is the product of many generations of selective breeding to bring out desirable qualities for tea. However, many other camellias can be used to produce a similar beverage. For example, tea made from C. sasanqua leaves is popular in some parts of Japan.
Seeds of C. oleifera, C. japonica, and, to a lesser extent, other species such as C. crapnelliana, C. reticulata, C. sasanqua and C. sinensis as well are pressed to make tea seed oil, a sweet seasoning and cooking oil special to East Asia. It is the most important cooking oil for hundreds of millions of people, particularly in southern China.
Camellia oil is commonly used to clean and protect the blades of cutting instruments.
Camellia oil pressed from seeds of C. japonica, also called tsubaki oil or tsubaki-abura (椿油) in Japanese, has been traditionally used in Japan for hair care. C. japonica plant is used to prepare traditional antiinflammatory medicines.
History
Fossil record
The earliest fossil record of Camellia are the leaves of †C. abensis from the upper Eocene
of Japan, †C. abchasica from the lower Oligocene of Bulgaria and
†C. multiforma from the lower Oligocene of Washington, United States.
Garden history
Camellias were cultivated in the gardens of China for centuries before they were seen in Europe. The German botanist Engelbert Kaempfer reported that the "Japan Rose", as he called it, grew wild in woodland and hedgerow, but that many superior varieties had been selected for gardens. Europeans' earliest views of camellias must have been their representations in Chinese painted wallpapers, where they were often represented growing in porcelain pots.
The first living camellias seen in England were a single red and a single white, grown and flowered in his garden at Thorndon Hall, Essex, by Robert James, Lord Petre, among the keenest gardeners of his generation, in 1739. His gardener James Gordon was the first to introduce camellias to commerce, from the nurseries he established after Lord Petre's untimely death in 1743, at Mile End, Essex, near London.
With the expansion of the tea trade in the later 18th century, new varieties began to be seen in England, imported through the British East India Company. The Company's John Slater was responsible for the first of the new camellias, double ones, in white and a striped red, imported in 1792. Further camellias imported in the East Indiamen were associated with the patrons whose gardeners grew them: a double red for Sir Robert Preston in 1794 and the pale pink named "Lady Hume's Blush" for Amelia, the lady of Sir Abraham Hume of Wormleybury, Hertfordshire (1806). The camellia was imported from England to America in 1797 when Colonel John Stevens brought the flower as part of an effort to grow attractions within Elysian Fields in Hoboken, New Jersey. By 1819, twenty-five camellias had bloomed in England; that year the first monograph appeared, Samuel Curtis's, A Monograph on the Genus Camellia, whose five handsome folio colored illustrations have usually been removed from the slender text and framed. Though they did not flower for over a decade, camellias that set seed rewarded their growers with a wealth of new varieties. By the 1840s, the camellia was at the height of its fashion as the luxury flower. The Parisian courtesan Marie Duplessis, who died young in 1847, inspired Dumas' La Dame aux camélias and Verdi's La Traviata.
The fashionable imbricated formality of prized camellias was an element in their decline, replaced by the new hothouse orchid. Their revival after World War I as woodland shrubs for mild climates has been paralleled by the rise in popularity of Camellia sasanqua.
Modern cultivars
The tea camellia, C. sinensis, has been selected by many commercial cultivars for the taste of its leaves once processed into tea leaves.
Today, camellias are grown as ornamental plants for their flowers; about 3,000 cultivars and hybrids have been selected, many with double or semi-double flowers. C. japonica is the most prominent species in cultivation, with over 2,000 named cultivars. Next are C. reticulata with over 400 named cultivars, and C. sasanqua with over 300 named cultivars. Popular hybrids include C. × hiemalis (C. japonica × C. sasanqua) and C. × williamsii (C. japonica × C. saluenensis). Some varieties can grow considerably, up to , though more compact cultivars are available. They are frequently planted in woodland settings alongside other calcifuges, such as rhododendrons. They are particularly associated with areas of high soil acidity, such as Cornwall and Devon in the UK. They are highly valued for their very early flowering, often among the first flowers to appear in the late winter. Late frosts can damage the flower buds, resulting in misshapen flowers.
There is a great variety of flower forms:
single (flat, bowl- or cup-shaped)
semi-double (rows of large outer petals, with the centre comprising mixed petals and stamens)
double:
paeony form (convex mass of irregular petals and petaloids with hidden stamens)
anemone form (one or more rows of outer petals, with mixed petaloids and stamens in the centre)
rose form (overlapping petals showing stamens in a concave centre when open)
formal double (rows of overlapping petals with hidden stamens)
AGM cultivars
The following hybrid cultivars have gained the Royal Horticultural Society's Award of Garden Merit:
Species
Plants of the World Online currently includes:
Camellia albata Orel & Curry
Camellia amplexicaulis (Pit.) Cohen-Stuart
Camellia amplexifolia Merr. & Chun
Camellia anlungensis Hung T.Chang
Camellia assimiloides Sealy
Camellia aurea Hung T.Chang
Camellia azalea C.F.Wei
Camellia brevistyla (Hayata) Cohen-Stuart
Camellia bugiamapensis Orel, Curry, Luu & Q.D.Nguyen
Camellia campanulata Orel, Curry & Luu
Camellia candida Hung T.Chang
Camellia capitata Orel, Curry & Luu
Camellia cattienensis Orel
Camellia caudata Wall.
Camellia chekiangoleosa Hu
Camellia cherryana Orel
Camellia chinmeiae S.L.Lee & T.Y.A.Yang
Camellia chrysanthoides Hung T.Chang
Camellia concinna Orel & Curry
Camellia connata (Craib) Craib
Camellia corallina (Gagnep.) Sealy
Camellia cordifolia (F.P.Metcalf) Nakai
Camellia costata S.Y.Hu & S.Y.Liang
Camellia costei H.Lév.
Camellia crapnelliana Tutcher – Crapnell's camellia
Camellia crassicolumna Hung T.Chang
Camellia crassipes Sealy
Camellia crassiphylla Ninh & Hakoda
Camellia cuongiana Orel & Curry
Camellia cupiformis T.L.Ming
Camellia curryana Orel & Luu
Camellia cuspidata (Kochs) Bean
Camellia dalatensis V.D.Luong, Ninh & Hakoda
Camellia debaoensis R.C.Hu & Y.Q.Liufu
Camellia decora Orel, Curry & Luu
Camellia dilinhensis Ninh & V.D.Luong
Camellia dongnaicensis Orel
Camellia dormoyana (Pierre ex Laness.) Sealy
Camellia drupifera Lour.
Camellia duyana Orel, Curry & Luu
Camellia edithae Hance
Camellia elizabethae Orel & Curry
Camellia elongata (Rehder & E.H.Wilson) Rehder
Camellia erubescens Orel & Curry
Camellia euphlebia Merr. ex Sealy
Camellia euryoides Lindl.
Camellia fangchengensis S.Ye Liang & Y.C.Zhong
Camellia fansipanensis J.M.H.Shaw, Wynn-Jones & V.D.Nguyen
Camellia fascicularis Hung T.Chang
Camellia flava (Pit.) Sealy
Camellia flavida Hung T.Chang
Camellia fleuryi (A.Chev.) Sealy
Camellia fluviatilis Hand.-Mazz.
Camellia forrestii (Diels) Cohen-Stuart
Camellia fraterna Hance
Camellia furfuracea (Merr.) Cohen-Stuart
Camellia gaudichaudii (Gagnep.) Sealy
Camellia gilbertii (A.Chev.) Sealy
Camellia glabricostata T.L.Ming
Camellia gracilipes Merr. ex Sealy
Camellia grandibracteata Hung T.Chang, Y.J.Tan, F.L.Yu & P.S.Wang
Camellia granthamiana Sealy – Grantham's camellia
Camellia grijsii Hance
Camellia gymnogyna Hung T.Chang
Camellia harlandii Orel & Curry
Camellia hatinhensis V.D.Luong, Ninh & L.T.Nguyen
Camellia hekouensis C.J.Wang & G.S.Fan
Camellia hiemalis Nakai
Camellia honbaensis Luu, Q.D.Nguyen & G.Tran
Camellia hongiaoensis Orel & Curry
Camellia hongkongensis Seem.
Camellia hsinpeiensis S.S.Ying
Camellia huana T.L.Ming & W.J.Zhang
Camellia ilicifolia Y.K.Li
Camellia impressinervis Hung T.Chang & S.Ye Liang
Camellia indochinensis Merr.
Camellia ingens Orel & Curry
Camellia insularis Orel & Curry
Camellia × intermedia (Tuyama) Nagam.
Camellia inusitata Orel, Curry & Luu
Camellia japonica L. – East Asian camelliasynonym Camellia rusticana – snow camellia
Camellia kissii Wall.
Camellia krempfii (Gagnep.) Sealy
Camellia kwangsiensis Hung T.Chang
Camellia lanceolata (Blume) Seem.
Camellia langbianensis (Gagnep.) P.H.Hô
Camellia laotica (Gagnep.) T.L.Ming
Camellia lawii Sealy
Camellia leptophylla S.Ye Liang ex Hung T.Chang
Camellia ligustrina Orel, Curry & Luu
Camellia longicalyx Hung T.Chang
Camellia longii Orel & Luu
Camellia longipedicellata (Hu) Hung T.Chang & D.Fang
Camellia longissima Hung T.Chang & S.Ye Liang
Camellia lucii Orel & Curry
Camellia lutchuensis T.Itô
Camellia luteocerata Orel
Camellia luteoflora Y.K.Li ex Hung T.Chang & F.A.Zeng
Camellia luteopallida V.D.Luong, T.Q.T.Nguyen & Luu
Camellia luuana Orel & Curry
Camellia maiana Orel
Camellia mairei (H.Lév.) Melch.
Camellia maoniushanensis J.L.Liu & Q.Luo
Camellia megasepala Hung T.Chang & Trin Ninh
Camellia melliana Hand.-Mazz.
Camellia micrantha S.Ye Liang & Y.C.Zhong
Camellia mileensis T.L.Ming
Camellia mingii S.X.Yang
Camellia minima Orel & Curry
Camellia mollis Hung T.Chang & S.X.Ren
Camellia montana (Blanco) Hung T.Chang & S.X.Ren
Camellia murauchii Ninh & Hakoda
Camellia namkadingensis Soulad. & Tagane
Camellia nematodea (Gagnep.) Sealy
Camellia nervosa (Gagnep.) Hung T.Chang
Camellia oconoriana Orel, Curry & Luu
Camellia oleifera C.Abel – oil-seed camellia, tea oil camellia
Camellia pachyandra Hu
Camellia parviflora Merr. & Chun ex Sealy
Camellia parvimuricata Hung T.Chang
Camellia paucipunctata (Merr. & Chun) Chun
Camellia petelotii (Merr.) Sealy synonyms:C. chrysantha, C. nitidissima – yellow camellia
Camellia philippinensis Hung T.Chang & S.X.Ren
Camellia pilosperma S.Yun Liang
Camellia pingguoensis D.Fang
Camellia piquetiana (Pierre) Sealy
Camellia pitardii Cohen-Stuart
Camellia pleurocarpa (Gagnep.) Sealy
Camellia polyodonta F.C.How ex Hu
Camellia psilocarpa X.G.Shi & C.X.Ye
Camellia ptilophylla Hung T.Chang
Camellia pubicosta Merr.
Camellia pubifurfuracea Y.C.Zhong
Camellia pubipetala Y.Wan & S.Z.Huang
Camellia pukhangensis N.D.Do, V.D.Luong, S.T.Hoang & T.H.Lê
Camellia punctata (Kochs) Cohen-Stuart
Camellia pyriparva Orel & Curry
Camellia pyxidiacea Z.R.Xu, F.P.Chen & C.Y.Deng
Camellia quangcuongii L.V.Dung, S.T. Hoang & Nhan
Camellia reflexa Orel & Curry
Camellia renshanxiangiae C.X.Ye & X.Q.Zheng
Camellia reticulata Lindl.
Camellia rhytidocarpa Hung T.Chang & S.Ye Liang
Camellia rosacea Tagane, Soulad. & Yahara
Camellia rosiflora Hook.
Camellia rosmannii Ninh
Camellia rosthorniana Hand.-Mazz.
Camellia rubriflora Ninh & Hakoda
Camellia salicifolia Champ.
Camellia saluenensis Stapf ex Bean
Camellia sasanqua Thunb.
Camellia scabrosa Orel & Curry
Camellia sealyana T.L.Ming
Camellia semiserrata C.W.Chi
Camellia septempetala Hung T.Chang & L.L.Qi
Camellia siangensis T.K.Paul & M.P.Nayar
Camellia sinensis (L.) Kuntze – tea plant
Camellia sonthaiensis Luu, V.D.Luong, Q.D.Nguyen & T.Q.T.Nguyen
Camellia stuartiana Sealy
Camellia subintegra P.C.Huang
Camellia synaptica Sealy
Camellia szechuanensis C.W.Chi
Camellia szemaoensis Hung T.Chang
Camellia tachangensis F.S.Zhang
Camellia tadungensis Orel, Curry & Luu
Camellia taliensis (W.W.Sm.) Melch. – also used to make tea like C. sinensis
Camellia tenii Sealy
Camellia thailandica Hung T.Chang & S.X.Ren
Camellia thanxaensa Hakoda & Kirino
Camellia tienyenensis Orel & Curry
Camellia tomentosa Orel & Curry
Camellia tonkinensis (Pit.) Cohen-Stuart
Camellia transarisanensis (Hayata) Cohen-Stuart
Camellia trichoclada (Rehder) S.S.Chien
Camellia tsaii Hu
Camellia tsingpienensis Hu
Camellia tuberculata S.S.Chien
Camellia tuyenquangensis V.D.Luong, Le & Ninh
Camellia uraku Kitam.
Camellia villicarpa S.S.Chien
Camellia viridicalyx Hung T.Chang & S.Ye Liang
Camellia viscosa Orel & Curry
Camellia vuquangensis V.D.Luong, Ninh & L.T.Nguyen
Camellia wardii Kobuski
Camellia xanthochroma K.M.Feng & L.S.Xie
Camellia yokdonensis Dung bis & Hakoda
Camellia yunkiangica Hung T.Chang, H.S.Wang & B.H.Chen
Camellia yunnanensis (Pit. ex Diels) Cohen-Stuart
Cultural significance
The Camellia family of plants in popular culture.
The following cities are nicknamed the "Camellia City" of each state: Greenville, Alabama; Sacramento, California; Fort Walton Beach, Florida; Slidell, Louisiana; McComb, Mississippi; and Newberg, Oregon. Meanwhile, Thomson, Georgia, is nicknamed the "Camellia City of the South".
The camellia is the state flower of Alabama.
A postseason college football bowl game introduced in 2014 in Montgomery, Alabama, was first known as the Camellia Bowl.
Alexandre Dumas fils wrote the novel and stage adaptation The Lady of the Camellias, wherein the flower is a symbol of a courtesan's sexual availability.
Augusta National Golf Club's 10th hole is named "Camellia", one of many references to the plant nursery originally on the site of the course.
Rabindranath Tagore wrote a poem titled "Camellia" about a youth's longing for a young woman he sees on the train.
In the book To Kill a Mockingbird, Jem destroys Mrs. Dubose's camellia bushes after she insults his family, yet he later receives a camellia bud from the dying woman.
A white camellia flower is an iconic symbol of Chanel haute couture, a tradition started by Coco Chanel herself, who identified with the heroine of Dumas's work.
Camellias have major significance in the Akira Kurosawa film Sanjuro, likely due to their association with the concept of "a noble death" in samurai culture.
White camellias are a symbol of the women's suffrage movement in New Zealand and appear on the country's ten-dollar note.
The Knights of the White Camelia was an organization similar to the Ku Klux Klan.
Temple City, California's slogan since 1944 has been "Temple City, Home of Camellias", and the city has become well-known for its Camellia Festival.
In Brazil, the camellia was a symbol of abolitionist movement during the Imperial Age. It was common practice for abolitionists to plant camellias in solidarity.
An Argentinian military march is called "Avenida de las Camelias".
Camellia flowers are featured on the cover of The Silent Circus, the second studio album by American progressive metal band Between the Buried and Me.
| Biology and health sciences | Ericales | null |
76142 | https://en.wikipedia.org/wiki/Hydrangea | Hydrangea | Hydrangea ( or ), commonly named the hortensia, is a genus of more than 70 species of flowering plants native to Asia and the Americas. By far the greatest species diversity is in eastern Asia, notably China, Korea, and Japan. Most are shrubs tall, but some are small trees, and others lianas reaching up to by climbing up trees. They can be either deciduous or evergreen, though the widely cultivated temperate species are all deciduous.
The flowers of many hydrangea act as natural pH indicators, sporting blue flowers when the soil is acidic and pink ones when the soil is alkaline.
Etymology
Hydrangea is derived from Greek and means ‘water vessel’ (from húdōr "water" + ángos or angeîon "vessel"), in reference to the shape of its seed capsules. The earlier name, Hortensia, is a Latinised version of the French given name Hortense, honoring French astronomer and mathematician Nicole-Reine Hortense Lepaute. Philibert Commerson attempted to name the flower Lepautia or Peautia after Lepaute. However, the flower's accepted name later became Hortensia. This led to people believing Lepaute's name was Hortense, but the Larousse remarks that this is erroneous, and that the name probably came from hortus, garden.
Life cycle
Hydrangea flowers are produced from early spring to late autumn; they grow in flowerheads (corymbs or panicles) most often at the ends of the stems. Typically the flowerheads contain two types of flowers: small non-showy fertile flowers in the center or interior of the flowerhead, and large, sterile showy flowers with large colorful sepals (tepals). These showy flowers are often extended in a ring, or to the exterior of the small flowers. Plants in wild populations typically have few to none of the showy flowers, while cultivated hydrangeas have been bred and selected to have more of the larger type flowers.
There are two flower arrangements in hydrangeas with corymb style inflorescences, which includes the commonly grown "bigleaf hydrangea"—Hydrangea macrophylla. Mophead flowers are large round flowerheads resembling pom-poms or, as the name implies, the head of a mop. In contrast, lacecap flowers bear round, flat flowerheads with a center core of subdued, small flowers surrounded by outer rings of larger flowers having showy sepals or tepals. The flowers of some rhododendrons and viburnums can appear, at first glance, similar to those of some hydrangeas.
Hydrangea flowers, when cut, dehydrate easily and wilt very quickly due to the large surface area of the petals. A wilted hydrangea may have its hydration restored by first having its stem immersed in boiling water; as the petals of the hydrangea can also absorb water, the petals may then be immersed, in room-temperature water, to restore the flower's hydration.
Colors and soil acidity
Hydrangea flower color can change based on the pH in soil. As the graph depicts, soil with a pH of 5.5 or lower will produce blue flowers, a pH of 6.5 or higher will produce pink hydrangeas, and soil in between 5.5 and 6.5 will have purple hydrangeas. White hydrangeas cannot be color-manipulated by soil pH because they do not produce pigment for color. In other words, while the hue of the inflorescence is variable dependent upon cultural factors, the color saturation is genetically predetermined. In most species, the flowers are white. In some, however, (notably H. macrophylla), they can be blue, red, or purple, with color saturation levels ranging from the palest of pinks, lavenders & powder blues, to deep, rich purples, reds, and royal blues. In these species, floral color change occurs due to the availability of aluminium ions, a variable which itself depends upon the soil pH. For H. macrophylla and H. serrata cultivars, the flower color can be determined by the relative acidity of the soil: an acidic soil (pH below 7), will have available aluminum ions and typically produce flowers that are blue to purple, whereas an alkaline soil (pH above 7) will tie up aluminium ions and result in pink or red flowers. This is caused by a color change of the flower pigments in the presence of aluminium ions which can be taken up into hyperaccumulating plants.
Partial list of species
Hydrangea anomala – (climbing hydrangea) Himalaya, southwest China
Hydrangea arborescens – (smooth hydrangea) eastern North America
Hydrangea aspera – China, Himalaya
Hydrangea bretschneideri – China
Hydrangea chinensis – China and Taiwan
Hydrangea chungii – China
Hydrangea cinerea – (ashy hydrangea) eastern United States
Hydrangea coenobialis – China
Hydrangea davidii – China
Hydrangea febrifuga – central & southern China to Malesia and New Guinea
Hydrangea glaucescens – China, Myanmar and Vietnam
Hydrangea gracilis – China
Hydrangea heteromalla – Himalaya, west and north China
Hydrangea hirta – Japan
Hydrangea hydrangeoides – Ulleungdo, Japan, Kurils
Hydrangea hypoglauca – China
Hydrangea integrifolia – China
Hydrangea involucrata – Japan, Taiwan
Hydrangea jelskii – Andes
Hydrangea kwangsiensis – China
Hydrangea kwangtungensis – China
Hydrangea lingii – China
Hydrangea linkweiensis – China
Hydrangea longifolia – China
Hydrangea longipes – western China
Hydrangea macrocarpa – China
Hydrangea macrophylla – (bigleaf hydrangea) southeast Japan, southern China
Hydrangea mangshanensis – China
Hydrangea paniculata – (panicled hydrangea) eastern China, Japan, Korea, Sakhalin
Hydrangea peruviana – Costa Rica and Panama, Andes
Hydrangea petiolaris – (climbing hydrangea) Japan, Korea, Sakhalin
Hydrangea quercifolia – (oakleaf hydrangea) southeast United States
Hydrangea radiata – (silverleaf hydrangea) southeast United States
Hydrangea robusta – China, Himalaya
Hydrangea sargentiana – western China
Hydrangea scandens – southern Japan south to the Philippines
Hydrangea serrata – Japan, Korea
Hydrangea serratifolia – Chile, western Argentina
Hydrangea strigosa – China
Hydrangea stylosa – China
Hydrangea tarapotensis – Andes
Hydrangea xanthoneura – China
Hydrangea zhewanensis – China
Fossil record
Hydrangea alaskana is a fossil species recovered from Paleogene strata at Jaw Mountain Alaska. †Hydrangea knowltoni has been described from leaves and flowers recovered from the Miocene Langhian Latah Formation of the inland Pacific Northwest United states. The related Miocene species †Hydrangea bendirei is known to from the Mascall Formation in Oregon, and †Hydrangea reticulata is documented from the Weaverville Formation in California.
Four fossil seeds of †Hydrangea polonica have been extracted from borehole samples of the Middle Miocene fresh water deposits in Nowy Sacz Basin, West Carpathians, Poland.
Cultivation and uses
Hydrangeas are popular ornamental plants, grown for their large flowerheads, with Hydrangea macrophylla being by far the most widely grown. It has over 600 named cultivars, many selected to have only large sterile flowers in the flowerheads. Hydrangea macrophylla, also known as bigleaf hydrangea, can be broken up into two main categories; mophead hydrangea and lacecap hydrangea. Some are best pruned on an annual basis when the new leaf buds begin to appear. If not pruned regularly, the bush will become very "leggy", growing upwards until the weight of the stems is greater than their strength, at which point the stems will sag down to the ground and possibly break. Other species only flower on "old wood". Thus, new wood resulting from pruning will not produce flowers until the following season.
The following cultivars and species have gained the Royal Horticultural Society's Award of Garden Merit under the synonym Schizophragma:
S. hydrangeoides var. concolor 'Moonlight'
S. hydrangeoides var. hydrangeoides 'Roseum'
S. integrifolium
Hydrangea root and rhizome are indicated for the treatment of conditions of the urinary tract in the Physicians' Desk Reference for Herbal Medicine and may have diuretic properties. Hydrangeas are moderately toxic if eaten, with all parts of the plant containing cyanogenic glycosides. Hydrangea paniculata is reportedly sometimes smoked as an intoxicant, despite the danger of illness and/or death due to the cyanide.
The flowers on a hydrangea shrub can change from blue to pink or from pink to blue from one season to the next depending on the acidity level of the soil. Adding organic materials such as coffee grounds and citrus peel will increase acidity and turn hydrangea flowers blue.
A popular pink hydrangea called Vanilla Strawberry has been named "Top Plant" by the American Nursery and Landscape Association.
A hybrid "Runaway Bride Snow White", from Japan, won Plant of the Year at the 2018 RHS Chelsea Flower Show.
In culture
In Japan, ama-cha (), meaning sweet tea, is another herbal tea made from Hydrangea serrata, whose leaves contain a substance that develops a sweet taste (phyllodulcin). For the fullest taste, fresh leaves are crumpled, steamed, and dried, yielding dark brown tea leaves. Ama-cha is mainly used for kan-butsu-e (the Buddha bathing ceremony) on April 8 every year—the day thought to be Buddha's birthday in Japan. During the ceremony, ama-cha is poured over a statue of Buddha and served to people in attendance. A legend has it that on the day Buddha was born, nine dragons poured Amrita over him; ama-cha is substituted for Amrita in Japan.
In Korean tea, Hydrangea serrata is used for an herbal tea called sugukcha () or isulcha ().
The pink hydrangea has risen in popularity all over the world, especially in Asia. The given meaning of pink hydrangeas is popularly tied to the phrase "you are the beat of my heart," as described by the celebrated Korean florist Tan Jun Yong, who was quoted saying, "The light delicate blush of the petals reminds me of a beating heart, while the size could only match the heart of the sender!"
Hydrangea quercifolia was declared the official state wildflower of the U.S. state of Alabama in 1999.
Hydrangeas were used by the Cherokee people of what is now the Southern U.S. A mild diuretic and cathartic, it was considered a valuable remedy for stone and gravel in the bladder.
Gallery
Diseases
| Biology and health sciences | Cornales | Plants |
76174 | https://en.wikipedia.org/wiki/Universal%20quantification | Universal quantification | In mathematical logic, a universal quantification is a type of quantifier, a logical constant which is interpreted as "given any", "for all", or "for any". It expresses that a predicate can be satisfied by every member of a domain of discourse. In other words, it is the predication of a property or relation to every member of the domain. It asserts that a predicate within the scope of a universal quantifier is true of every value of a predicate variable.
It is usually denoted by the turned A (∀) logical operator symbol, which, when used together with a predicate variable, is called a universal quantifier ("", "", or sometimes by "" alone). Universal quantification is distinct from existential quantification ("there exists"), which only asserts that the property or relation holds for at least one member of the domain.
Quantification in general is covered in the article on quantification (logic). The universal quantifier is encoded as in Unicode, and as \forall in LaTeX and related formula editors.
Basics
Suppose it is given that
2·0 = 0 + 0, and 2·1 = 1 + 1, and , ..., and 2 · 100 = 100 + 100, and ..., etc.
This would seem to be an infinite logical conjunction because of the repeated use of "and". However, the "etc." cannot be interpreted as a conjunction in formal logic, Instead, the statement must be rephrased:
For all natural numbers n, one has 2·n = n + n.
This is a single statement using universal quantification.
This statement can be said to be more precise than the original one. While the "etc." informally includes natural numbers, and nothing more, this was not rigorously given. In the universal quantification, on the other hand, the natural numbers are mentioned explicitly.
This particular example is true, because any natural number could be substituted for n and the statement "2·n = n + n" would be true. In contrast,
For all natural numbers n, one has 2·n > 2 + n
is false, because if n is substituted with, for instance, 1, the statement "2·1 > 2 + 1" is false. It is immaterial that "2·n > 2 + n" is true for most natural numbers n: even the existence of a single counterexample is enough to prove the universal quantification false.
On the other hand,
for all composite numbers n, one has 2·n > 2 + n
is true, because none of the counterexamples are composite numbers. This indicates the importance of the domain of discourse, which specifies which values n can take. In particular, note that if the domain of discourse is restricted to consist only of those objects that satisfy a certain predicate, then for universal quantification this requires a logical conditional. For example,
For all composite numbers n, one has 2·n > 2 + n
is logically equivalent to
For all natural numbers n, if n is composite, then 2·n > 2 + n.
Here the "if ... then" construction indicates the logical conditional.
Notation
In symbolic logic, the universal quantifier symbol (a turned "A" in a sans-serif font, Unicode U+2200) is used to indicate universal quantification. It was first used in this way by Gerhard Gentzen in 1935, by analogy with Giuseppe Peano's (turned E) notation for existential quantification and the later use of Peano's notation by Bertrand Russell.
For example, if P(n) is the predicate "2·n > 2 + n" and N is the set of natural numbers, then
is the (false) statement
"for all natural numbers n, one has 2·n > 2 + n".
Similarly, if Q(n) is the predicate "n is composite", then
is the (true) statement
"for all natural numbers n, if n is composite, then ".
Several variations in the notation for quantification (which apply to all forms) can be found in the Quantifier article.
Properties
Negation
The negation of a universally quantified function is obtained by changing the universal quantifier into an existential quantifier and negating the quantified formula. That is,
where denotes negation.
For example, if is the propositional function " is married", then, for the set of all living human beings, the universal quantification
Given any living person , that person is married
is written
This statement is false. Truthfully, it is stated that
It is not the case that, given any living person , that person is married
or, symbolically:
.
If the function is not true for every element of , then there must be at least one element for which the statement is false. That is, the negation of is logically equivalent to "There exists a living person who is not married", or:
It is erroneous to confuse "all persons are not married" (i.e. "there exists no person who is married") with "not all persons are married" (i.e. "there exists a person who is not married"):
Other connectives
The universal (and existential) quantifier moves unchanged across the logical connectives ∧, ∨, →, and ↚, as long as the other operand is not affected; that is:
Conversely, for the logical connectives ↑, ↓, ↛, and ←, the quantifiers flip:
Rules of inference
A rule of inference is a rule justifying a logical step from hypothesis to conclusion. There are several rules of inference which utilize the universal quantifier.
Universal instantiation concludes that, if the propositional function is known to be universally true, then it must be true for any arbitrary element of the universe of discourse. Symbolically, this is represented as
where c is a completely arbitrary element of the universe of discourse.
Universal generalization concludes the propositional function must be universally true if it is true for any arbitrary element of the universe of discourse. Symbolically, for an arbitrary c,
The element c must be completely arbitrary; else, the logic does not follow: if c is not arbitrary, and is instead a specific element of the universe of discourse, then P(c) only implies an existential quantification of the propositional function.
The empty set
By convention, the formula is always true, regardless of the formula P(x); see vacuous truth.
Universal closure
The universal closure of a formula φ is the formula with no free variables obtained by adding a universal quantifier for every free variable in φ. For example, the universal closure of
is
.
As adjoint
In category theory and the theory of elementary topoi, the universal quantifier can be understood as the right adjoint of a functor between power sets, the inverse image functor of a function between sets; likewise, the existential quantifier is the left adjoint.
For a set , let denote its powerset. For any function between sets and , there is an inverse image functor between powersets, that takes subsets of the codomain of f back to subsets of its domain. The left adjoint of this functor is the existential quantifier and the right adjoint is the universal quantifier .
That is, is a functor that, for each subset , gives the subset given by
those in the image of under . Similarly, the universal quantifier is a functor that, for each subset , gives the subset given by
those whose preimage under is contained in .
The more familiar form of the quantifiers as used in first-order logic is obtained by taking the function f to be the unique function so that is the two-element set holding the values true and false, a subset S is that subset for which the predicate holds, and
which is true if is not empty, and
which is false if S is not X.
The universal and existential quantifiers given above generalize to the presheaf category.
| Mathematics | Mathematical logic | null |
76340 | https://en.wikipedia.org/wiki/Principal%20component%20analysis | Principal component analysis | Principal component analysis (PCA) is a linear dimensionality reduction technique with applications in exploratory data analysis, visualization and data preprocessing.
The data is linearly transformed onto a new coordinate system such that the directions (principal components) capturing the largest variation in the data can be easily identified.
The principal components of a collection of points in a real coordinate space are a sequence of unit vectors, where the -th vector is the direction of a line that best fits the data while being orthogonal to the first vectors. Here, a best-fitting line is defined as one that minimizes the average squared perpendicular distance from the points to the line. These directions (i.e., principal components) constitute an orthonormal basis in which different individual dimensions of the data are linearly uncorrelated. Many studies use the first two principal components in order to plot the data in two dimensions and to visually identify clusters of closely related data points.
Principal component analysis has applications in many fields such as population genetics, microbiome studies, and atmospheric science.
Overview
When performing PCA, the first principal component of a set of variables is the derived variable formed as a linear combination of the original variables that explains the most variance. The second principal component explains the most variance in what is left once the effect of the first component is removed, and we may proceed through iterations until all the variance is explained. PCA is most commonly used when many of the variables are highly correlated with each other and it is desirable to reduce their number to an independent set.
The first principal component can equivalently be defined as a direction that maximizes the variance of the projected data. The -th principal component can be taken as a direction orthogonal to the first principal components that maximizes the variance of the projected data.
For either objective, it can be shown that the principal components are eigenvectors of the data's covariance matrix. Thus, the principal components are often computed by eigendecomposition of the data covariance matrix or singular value decomposition of the data matrix. PCA is the simplest of the true eigenvector-based multivariate analyses and is closely related to factor analysis. Factor analysis typically incorporates more domain-specific assumptions about the underlying structure and solves eigenvectors of a slightly different matrix. PCA is also related to canonical correlation analysis (CCA). CCA defines coordinate systems that optimally describe the cross-covariance between two datasets while PCA defines a new orthogonal coordinate system that optimally describes variance in a single dataset. Robust and L1-norm-based variants of standard PCA have also been proposed.
History
PCA was invented in 1901 by Karl Pearson, as an analogue of the principal axis theorem in mechanics; it was later independently developed and named by Harold Hotelling in the 1930s. Depending on the field of application, it is also named the discrete Karhunen–Loève transform (KLT) in signal processing, the Hotelling transform in multivariate quality control, proper orthogonal decomposition (POD) in mechanical engineering, singular value decomposition (SVD) of X (invented in the last quarter of the 19th century), eigenvalue decomposition (EVD) of XTX in linear algebra, factor analysis (for a discussion of the differences between PCA and factor analysis see Ch. 7 of Jolliffe's Principal Component Analysis), Eckart–Young theorem (Harman, 1960), or empirical orthogonal functions (EOF) in meteorological science (Lorenz, 1956), empirical eigenfunction decomposition (Sirovich, 1987), quasiharmonic modes (Brooks et al., 1988), spectral decomposition in noise and vibration, and empirical modal analysis in structural dynamics.
Intuition
PCA can be thought of as fitting a p-dimensional ellipsoid to the data, where each axis of the ellipsoid represents a principal component. If some axis of the ellipsoid is small, then the variance along that axis is also small.
To find the axes of the ellipsoid, we must first center the values of each variable in the dataset on 0 by subtracting the mean of the variable's observed values from each of those values. These transformed values are used instead of the original observed values for each of the variables. Then, we compute the covariance matrix of the data and calculate the eigenvalues and corresponding eigenvectors of this covariance matrix. Then we must normalize each of the orthogonal eigenvectors to turn them into unit vectors. Once this is done, each of the mutually-orthogonal unit eigenvectors can be interpreted as an axis of the ellipsoid fitted to the data. This choice of basis will transform the covariance matrix into a diagonalized form, in which the diagonal elements represent the variance of each axis. The proportion of the variance that each eigenvector represents can be calculated by dividing the eigenvalue corresponding to that eigenvector by the sum of all eigenvalues.
Biplots and scree plots (degree of explained variance) are used to interpret findings of the PCA.
Details
PCA is defined as an orthogonal linear transformation on a real inner product space that transforms the data to a new coordinate system such that the greatest variance by some scalar projection of the data comes to lie on the first coordinate (called the first principal component), the second greatest variance on the second coordinate, and so on.
Consider an data matrix, X, with column-wise zero empirical mean (the sample mean of each column has been shifted to zero), where each of the n rows represents a different repetition of the experiment, and each of the p columns gives a particular kind of feature (say, the results from a particular sensor).
Mathematically, the transformation is defined by a set of size of p-dimensional vectors of weights or coefficients that map each row vector of X to a new vector of principal component scores , given by
in such a way that the individual variables of t considered over the data set successively inherit the maximum possible variance from X, with each coefficient vector w constrained to be a unit vector (where is usually selected to be strictly less than to reduce dimensionality).
The above may equivalently be written in matrix form as
where
,
, and
.
First component
In order to maximize variance, the first weight vector w(1) thus has to satisfy
Equivalently, writing this in matrix form gives
Since w(1) has been defined to be a unit vector, it equivalently also satisfies
The quantity to be maximised can be recognised as a Rayleigh quotient. A standard result for a positive semidefinite matrix such as XTX is that the quotient's maximum possible value is the largest eigenvalue of the matrix, which occurs when w is the corresponding eigenvector.
With w(1) found, the first principal component of a data vector x(i) can then be given as a score t1(i) = x(i) ⋅ w(1) in the transformed co-ordinates, or as the corresponding vector in the original variables, {x(i) ⋅ w(1)} w(1).
Further components
The k-th component can be found by subtracting the first k − 1 principal components from X:
and then finding the weight vector which extracts the maximum variance from this new data matrix
It turns out that this gives the remaining eigenvectors of XTX, with the maximum values for the quantity in brackets given by their corresponding eigenvalues. Thus the weight vectors are eigenvectors of XTX.
The k-th principal component of a data vector x(i) can therefore be given as a score tk(i) = x(i) ⋅ w(k) in the transformed coordinates, or as the corresponding vector in the space of the original variables, {x(i) ⋅ w(k)} w(k), where w(k) is the kth eigenvector of XTX.
The full principal components decomposition of X can therefore be given as
where W is a p-by-p matrix of weights whose columns are the eigenvectors of XTX. The transpose of W is sometimes called the whitening or sphering transformation. Columns of W multiplied by the square root of corresponding eigenvalues, that is, eigenvectors scaled up by the variances, are called loadings in PCA or in Factor analysis.
Covariances
XTX itself can be recognized as proportional to the empirical sample covariance matrix of the dataset XT.
The sample covariance Q between two of the different principal components over the dataset is given by:
where the eigenvalue property of w(k) has been used to move from line 2 to line 3. However eigenvectors w(j) and w(k) corresponding to eigenvalues of a symmetric matrix are orthogonal (if the eigenvalues are different), or can be orthogonalised (if the vectors happen to share an equal repeated value). The product in the final line is therefore zero; there is no sample covariance between different principal components over the dataset.
Another way to characterise the principal components transformation is therefore as the transformation to coordinates which diagonalise the empirical sample covariance matrix.
In matrix form, the empirical covariance matrix for the original variables can be written
The empirical covariance matrix between the principal components becomes
where Λ is the diagonal matrix of eigenvalues λ(k) of XTX. λ(k) is equal to the sum of the squares over the dataset associated with each component k, that is, λ(k) = Σi tk2(i) = Σi (x(i) ⋅ w(k))2.
Dimensionality reduction
The transformation T = X W maps a data vector x(i) from an original space of p variables to a new space of p variables which are uncorrelated over the dataset. However, not all the principal components need to be kept. Keeping only the first L principal components, produced by using only the first L eigenvectors, gives the truncated transformation
where the matrix TL now has n rows but only L columns. In other words, PCA learns a linear transformation where the columns of matrix form an orthogonal basis for the L features (the components of representation t) that are decorrelated. By construction, of all the transformed data matrices with only L columns, this score matrix maximises the variance in the original data that has been preserved, while minimising the total squared reconstruction error or .
Such dimensionality reduction can be a very useful step for visualising and processing high-dimensional datasets, while still retaining as much of the variance in the dataset as possible. For example, selecting L = 2 and keeping only the first two principal components finds the two-dimensional plane through the high-dimensional dataset in which the data is most spread out, so if the data contains clusters these too may be most spread out, and therefore most visible to be plotted out in a two-dimensional diagram; whereas if two directions through the data (or two of the original variables) are chosen at random, the clusters may be much less spread apart from each other, and may in fact be much more likely to substantially overlay each other, making them indistinguishable.
Similarly, in regression analysis, the larger the number of explanatory variables allowed, the greater is the chance of overfitting the model, producing conclusions that fail to generalise to other datasets. One approach, especially when there are strong correlations between different possible explanatory variables, is to reduce them to a few principal components and then run the regression against them, a method called principal component regression.
Dimensionality reduction may also be appropriate when the variables in a dataset are noisy. If each column of the dataset contains independent identically distributed Gaussian noise, then the columns of T will also contain similarly identically distributed Gaussian noise (such a distribution is invariant under the effects of the matrix W, which can be thought of as a high-dimensional rotation of the co-ordinate axes). However, with more of the total variance concentrated in the first few principal components compared to the same noise variance, the proportionate effect of the noise is less—the first few components achieve a higher signal-to-noise ratio. PCA thus can have the effect of concentrating much of the signal into the first few principal components, which can usefully be captured by dimensionality reduction; while the later principal components may be dominated by noise, and so disposed of without great loss. If the dataset is not too large, the significance of the principal components can be tested using parametric bootstrap, as an aid in determining how many principal components to retain.
Singular value decomposition
The principal components transformation can also be associated with another matrix factorization, the singular value decomposition (SVD) of X,
Here Σ is an n-by-p rectangular diagonal matrix of positive numbers σ(k), called the singular values of X; U is an n-by-n matrix, the columns of which are orthogonal unit vectors of length n called the left singular vectors of X; and W is a p-by-p matrix whose columns are orthogonal unit vectors of length p and called the right singular vectors of X.
In terms of this factorization, the matrix XTX can be written
where is the square diagonal matrix with the singular values of X and the excess zeros chopped off that satisfies . Comparison with the eigenvector factorization of XTX establishes that the right singular vectors W of X are equivalent to the eigenvectors of XTX, while the singular values σ(k) of are equal to the square-root of the eigenvalues λ(k) of XTX.
Using the singular value decomposition the score matrix T can be written
so each column of T is given by one of the left singular vectors of X multiplied by the corresponding singular value. This form is also the polar decomposition of T.
Efficient algorithms exist to calculate the SVD of X without having to form the matrix XTX, so computing the SVD is now the standard way to calculate a principal components analysis from a data matrix, unless only a handful of components are required.
As with the eigen-decomposition, a truncated score matrix TL can be obtained by considering only the first L largest singular values and their singular vectors:
The truncation of a matrix M or T using a truncated singular value decomposition in this way produces a truncated matrix that is the nearest possible matrix of rank L to the original matrix, in the sense of the difference between the two having the smallest possible Frobenius norm, a result known as the Eckart–Young theorem [1936].
Further considerations
The singular values (in Σ) are the square roots of the eigenvalues of the matrix XTX. Each eigenvalue is proportional to the portion of the "variance" (more correctly of the sum of the squared distances of the points from their multidimensional mean) that is associated with each eigenvector. The sum of all the eigenvalues is equal to the sum of the squared distances of the points from their multidimensional mean. PCA essentially rotates the set of points around their mean in order to align with the principal components. This moves as much of the variance as possible (using an orthogonal transformation) into the first few dimensions. The values in the remaining dimensions, therefore, tend to be small and may be dropped with minimal loss of information (see below). PCA is often used in this manner for dimensionality reduction. PCA has the distinction of being the optimal orthogonal transformation for keeping the subspace that has largest "variance" (as defined above). This advantage, however, comes at the price of greater computational requirements if compared, for example, and when applicable, to the discrete cosine transform, and in particular to the DCT-II which is simply known as the "DCT". Nonlinear dimensionality reduction techniques tend to be more computationally demanding than PCA.
PCA is sensitive to the scaling of the variables. If we have just two variables and they have the same sample variance and are completely correlated, then the PCA will entail a rotation by 45° and the "weights" (they are the cosines of rotation) for the two variables with respect to the principal component will be equal. But if we multiply all values of the first variable by 100, then the first principal component will be almost the same as that variable, with a small contribution from the other variable, whereas the second component will be almost aligned with the second original variable. This means that whenever the different variables have different units (like temperature and mass), PCA is a somewhat arbitrary method of analysis. (Different results would be obtained if one used Fahrenheit rather than Celsius for example.) Pearson's original paper was entitled "On Lines and Planes of Closest Fit to Systems of Points in Space" – "in space" implies physical Euclidean space where such concerns do not arise. One way of making the PCA less arbitrary is to use variables scaled so as to have unit variance, by standardizing the data and hence use the autocorrelation matrix instead of the autocovariance matrix as a basis for PCA. However, this compresses (or expands) the fluctuations in all dimensions of the signal space to unit variance.
Mean subtraction (a.k.a. "mean centering") is necessary for performing classical PCA to ensure that the first principal component describes the direction of maximum variance. If mean subtraction is not performed, the first principal component might instead correspond more or less to the mean of the data. A mean of zero is needed for finding a basis that minimizes the mean square error of the approximation of the data.
Mean-centering is unnecessary if performing a principal components analysis on a correlation matrix, as the data are already centered after calculating correlations. Correlations are derived from the cross-product of two standard scores (Z-scores) or statistical moments (hence the name: Pearson Product-Moment Correlation). Also see the article by Kromrey & Foster-Johnson (1998) on "Mean-centering in Moderated Regression: Much Ado About Nothing". Since covariances are correlations of normalized variables (Z- or standard-scores) a PCA based on the correlation matrix of X is equal to a PCA based on the covariance matrix of Z, the standardized version of X.
PCA is a popular primary technique in pattern recognition. It is not, however, optimized for class separability. However, it has been used to quantify the distance between two or more classes by calculating center of mass for each class in principal component space and reporting Euclidean distance between center of mass of two or more classes. The linear discriminant analysis is an alternative which is optimized for class separability.
Table of symbols and abbreviations
Properties and limitations
Properties
Some properties of PCA include:
Property 1: For any integer q, 1 ≤ q ≤ p, consider the orthogonal linear transformation
where is a q-element vector and is a (q × p) matrix, and let be the variance-covariance matrix for . Then the trace of , denoted , is maximized by taking , where consists of the first q columns of is the transpose of . ( is not defined here)
Property 2: Consider again the orthonormal transformation
with and defined as before. Then is minimized by taking where consists of the last q columns of .
The statistical implication of this property is that the last few PCs are not simply unstructured left-overs after removing the important PCs. Because these last PCs have variances as small as possible they are useful in their own right. They can help to detect unsuspected near-constant linear relationships between the elements of , and they may also be useful in regression, in selecting a subset of variables from , and in outlier detection.
Property 3: (Spectral decomposition of )
Before we look at its usage, we first look at diagonal elements,
Then, perhaps the main statistical implication of the result is that not only can we decompose the combined variances of all the elements of into decreasing contributions due to each PC, but we can also decompose the whole covariance matrix into contributions from each PC. Although not strictly decreasing, the elements of will tend to become smaller as increases, as is nonincreasing for increasing , whereas the elements of tend to stay about the same size because of the normalization constraints: .
Limitations
As noted above, the results of PCA depend on the scaling of the variables. This can be cured by scaling each feature by its standard deviation, so that one ends up with dimensionless features with unital variance.
The applicability of PCA as described above is limited by certain (tacit) assumptions made in its derivation. In particular, PCA can capture linear correlations between the features but fails when this assumption is violated (see Figure 6a in the reference). In some cases, coordinate transformations can restore the linearity assumption and PCA can then be applied (see kernel PCA).
Another limitation is the mean-removal process before constructing the covariance matrix for PCA. In fields such as astronomy, all the signals are non-negative, and the mean-removal process will force the mean of some astrophysical exposures to be zero, which consequently creates unphysical negative fluxes, and forward modeling has to be performed to recover the true magnitude of the signals. As an alternative method, non-negative matrix factorization focusing only on the non-negative elements in the matrices, which is well-suited for astrophysical observations. See more at Relation between PCA and Non-negative Matrix Factorization.
PCA is at a disadvantage if the data has not been standardized before applying the algorithm to it. PCA transforms original data into data that is relevant to the principal components of that data, which means that the new data variables cannot be interpreted in the same ways that the originals were. They are linear interpretations of the original variables. Also, if PCA is not performed properly, there is a high likelihood of information loss.
PCA relies on a linear model. If a dataset has a pattern hidden inside it that is nonlinear, then PCA can actually steer the analysis in the complete opposite direction of progress. Researchers at Kansas State University discovered that the sampling error in their experiments impacted the bias of PCA results. "If the number of subjects or blocks is smaller than 30, and/or the researcher is interested in PC's beyond the first, it may be better to first correct for the serial correlation, before PCA is conducted". The researchers at Kansas State also found that PCA could be "seriously biased if the autocorrelation structure of the data is not correctly handled".
PCA and information theory
Dimensionality reduction results in a loss of information, in general. PCA-based dimensionality reduction tends to minimize that information loss, under certain signal and noise models.
Under the assumption that
that is, that the data vector is the sum of the desired information-bearing signal and a noise signal one can show that PCA can be optimal for dimensionality reduction, from an information-theoretic point-of-view.
In particular, Linsker showed that if is Gaussian and is Gaussian noise with a covariance matrix proportional to the identity matrix, the PCA maximizes the mutual information between the desired information and the dimensionality-reduced output .
If the noise is still Gaussian and has a covariance matrix proportional to the identity matrix (that is, the components of the vector are iid), but the information-bearing signal is non-Gaussian (which is a common scenario), PCA at least minimizes an upper bound on the information loss, which is defined as
The optimality of PCA is also preserved if the noise is iid and at least more Gaussian (in terms of the Kullback–Leibler divergence) than the information-bearing signal . In general, even if the above signal model holds, PCA loses its information-theoretic optimality as soon as the noise becomes dependent.
Computation using the covariance method
The following is a detailed description of PCA using the covariance method as opposed to the correlation method.
The goal is to transform a given data set X of dimension p to an alternative data set Y of smaller dimension L. Equivalently, we are seeking to find the matrix Y, where Y is the Karhunen–Loève transform (KLT) of matrix X:
Organize the data set
Suppose you have data comprising a set of observations of p variables, and you want to reduce the data so that each observation can be described with only L variables, L < p. Suppose further, that the data are arranged as a set of n data vectors with each representing a single grouped observation of the p variables.
Write as row vectors, each with p elements.
Place the row vectors into a single matrix X of dimensions n × p.
Calculate the empirical mean
Find the empirical mean along each column j = 1, ..., p.
Place the calculated mean values into an empirical mean vector u of dimensions p × 1.
Calculate the deviations from the mean
Mean subtraction is an integral part of the solution towards finding a principal component basis that minimizes the mean square error of approximating the data. Hence we proceed by centering the data as follows:
Subtract the empirical mean vector from each row of the data matrix X.
Store mean-subtracted data in the n × p matrix B. where h is an column vector of all 1s:
In some applications, each variable (column of B) may also be scaled to have a variance equal to 1 (see Z-score). This step affects the calculated principal components, but makes them independent of the units used to measure the different variables.
Find the covariance matrix
Find the p × p empirical covariance matrix C from matrix B: where is the conjugate transpose operator. If B consists entirely of real numbers, which is the case in many applications, the "conjugate transpose" is the same as the regular transpose.
The reasoning behind using instead of n to calculate the covariance is Bessel's correction.
Find the eigenvectors and eigenvalues of the covariance matrix
Compute the matrix V of eigenvectors which diagonalizes the covariance matrix C: where D is the diagonal matrix of eigenvalues of C. This step will typically involve the use of a computer-based algorithm for computing eigenvectors and eigenvalues. These algorithms are readily available as sub-components of most matrix algebra systems, such as SAS, R, MATLAB, Mathematica, SciPy, IDL (Interactive Data Language), or GNU Octave as well as OpenCV.
Matrix D will take the form of an p × p diagonal matrix, where is the jth eigenvalue of the covariance matrix C, and
Matrix V, also of dimension p × p, contains p column vectors, each of length p, which represent the p eigenvectors of the covariance matrix C.
The eigenvalues and eigenvectors are ordered and paired. The jth eigenvalue corresponds to the jth eigenvector.
Matrix V denotes the matrix of right eigenvectors (as opposed to left eigenvectors). In general, the matrix of right eigenvectors need not be the (conjugate) transpose of the matrix of left eigenvectors.
Rearrange the eigenvectors and eigenvalues
Sort the columns of the eigenvector matrix V and eigenvalue matrix D in order of decreasing eigenvalue.
Make sure to maintain the correct pairings between the columns in each matrix.
Compute the cumulative energy content for each eigenvector
The eigenvalues represent the distribution of the source data's energy among each of the eigenvectors, where the eigenvectors form a basis for the data. The cumulative energy content g for the jth eigenvector is the sum of the energy content across all of the eigenvalues from 1 through j:
Select a subset of the eigenvectors as basis vectors
Save the first L columns of V as the p × L matrix W: where
Use the vector g as a guide in choosing an appropriate value for L. The goal is to choose a value of L as small as possible while achieving a reasonably high value of g on a percentage basis. For example, you may want to choose L so that the cumulative energy g is above a certain threshold, like 90 percent. In this case, choose the smallest value of L such that
Project the data onto the new basis
The projected data points are the rows of the matrix
That is, the first column of is the projection of the data points onto the first principal component, the second column is the projection onto the second principal component, etc.
Derivation using the covariance method
Let X be a d-dimensional random vector expressed as column vector. Without loss of generality, assume X has zero mean.
We want to find a orthonormal transformation matrix P so that PX has a diagonal covariance matrix (that is, PX is a random vector with all its distinct components pairwise uncorrelated).
A quick computation assuming were unitary yields:
Hence holds if and only if were diagonalisable by .
This is very constructive, as cov(X) is guaranteed to be a non-negative definite matrix and thus is guaranteed to be diagonalisable by some unitary matrix.
Covariance-free computation
In practical implementations, especially with high dimensional data (large ), the naive covariance method is rarely used because it is not efficient due to high computational and memory costs of explicitly determining the covariance matrix. The covariance-free approach avoids the operations of explicitly calculating and storing the covariance matrix , instead utilizing one of matrix-free methods, for example, based on the function evaluating the product at the cost of operations.
Iterative computation
One way to compute the first principal component efficiently is shown in the following pseudo-code, for a data matrix with zero mean, without ever computing its covariance matrix.
= a random vector of length
r = r / norm(r)
do times:
(a vector of length )
return
This power iteration algorithm simply calculates the vector , normalizes, and places the result back in . The eigenvalue is approximated by , which is the Rayleigh quotient on the unit vector for the covariance matrix . If the largest singular value is well separated from the next largest one, the vector gets close to the first principal component of within the number of iterations , which is small relative to , at the total cost . The power iteration convergence can be accelerated without noticeably sacrificing the small cost per iteration using more advanced matrix-free methods, such as the Lanczos algorithm or the Locally Optimal Block Preconditioned Conjugate Gradient (LOBPCG) method.
Subsequent principal components can be computed one-by-one via deflation or simultaneously as a block. In the former approach, imprecisions in already computed approximate principal components additively affect the accuracy of the subsequently computed principal components, thus increasing the error with every new computation. The latter approach in the block power method replaces single-vectors and with block-vectors, matrices and . Every column of approximates one of the leading principal components, while all columns are iterated simultaneously. The main calculation is evaluation of the product . Implemented, for example, in LOBPCG, efficient blocking eliminates the accumulation of the errors, allows using high-level BLAS matrix-matrix product functions, and typically leads to faster convergence, compared to the single-vector one-by-one technique.
The NIPALS method
Non-linear iterative partial least squares (NIPALS) is a variant the classical power iteration with matrix deflation by subtraction implemented for computing the first few components in a principal component or partial least squares analysis. For very-high-dimensional datasets, such as those generated in the *omics sciences (for example, genomics, metabolomics) it is usually only necessary to compute the first few PCs. The non-linear iterative partial least squares (NIPALS) algorithm updates iterative approximations to the leading scores and loadings t1 and r1T by the power iteration multiplying on every iteration by X on the left and on the right, that is, calculation of the covariance matrix is avoided, just as in the matrix-free implementation of the power iterations to , based on the function evaluating the product .
The matrix deflation by subtraction is performed by subtracting the outer product, t1r1T from X leaving the deflated residual matrix used to calculate the subsequent leading PCs.
For large data matrices, or matrices that have a high degree of column collinearity, NIPALS suffers from loss of orthogonality of PCs due to machine precision round-off errors accumulated in each iteration and matrix deflation by subtraction. A Gram–Schmidt re-orthogonalization algorithm is applied to both the scores and the loadings at each iteration step to eliminate this loss of orthogonality. NIPALS reliance on single-vector multiplications cannot take advantage of high-level BLAS and results in slow convergence for clustered leading singular values—both these deficiencies are resolved in more sophisticated matrix-free block solvers, such as the Locally Optimal Block Preconditioned Conjugate Gradient (LOBPCG) method.
Online/sequential estimation
In an "online" or "streaming" situation with data arriving piece by piece rather than being stored in a single batch, it is useful to make an estimate of the PCA projection that can be updated sequentially. This can be done efficiently, but requires different algorithms.
Qualitative variables
In PCA, it is common that we want to introduce qualitative variables as supplementary elements. For example, many quantitative variables have been measured on plants. For these plants, some qualitative variables are available as, for example, the species to which the plant belongs. These data were subjected to PCA for quantitative variables. When analyzing the results, it is natural to connect the principal components to the qualitative variable species.
For this, the following results are produced.
Identification, on the factorial planes, of the different species, for example, using different colors.
Representation, on the factorial planes, of the centers of gravity of plants belonging to the same species.
For each center of gravity and each axis, p-value to judge the significance of the difference between the center of gravity and origin.
These results are what is called introducing a qualitative variable as supplementary element. This procedure is detailed in and Husson, Lê, & Pagès (2009) and Pagès (2013).
Few software offer this option in an "automatic" way. This is the case of SPAD that historically, following the work of Ludovic Lebart, was the first to propose this option, and the R package FactoMineR.
Applications
Intelligence
The earliest application of factor analysis was in locating and measuring components of human intelligence. It was believed that intelligence had various uncorrelated components such as spatial intelligence, verbal intelligence, induction, deduction etc and that scores on these could be adduced by factor analysis from results on various tests, to give a single index known as the Intelligence Quotient (IQ). The pioneering statistical psychologist Spearman actually developed factor analysis in 1904 for his two-factor theory of intelligence, adding a formal technique to the science of psychometrics. In 1924 Thurstone looked for 56 factors of intelligence, developing the notion of Mental Age. Standard IQ tests today are based on this early work.
Residential differentiation
In 1949, Shevky and Williams introduced the theory of factorial ecology, which dominated studies of residential differentiation from the 1950s to the 1970s. Neighbourhoods in a city were recognizable or could be distinguished from one another by various characteristics which could be reduced to three by factor analysis. These were known as 'social rank' (an index of occupational status), 'familism' or family size, and 'ethnicity'; Cluster analysis could then be applied to divide the city into clusters or precincts according to values of the three key factor variables. An extensive literature developed around factorial ecology in urban geography, but the approach went out of fashion after 1980 as being methodologically primitive and having little place in postmodern geographical paradigms.
One of the problems with factor analysis has always been finding convincing names for the various artificial factors. In 2000, Flood revived the factorial ecology approach to show that principal components analysis actually gave meaningful answers directly, without resorting to factor rotation. The principal components were actually dual variables or shadow prices of 'forces' pushing people together or apart in cities. The first component was 'accessibility', the classic trade-off between demand for travel and demand for space, around which classical urban economics is based. The next two components were 'disadvantage', which keeps people of similar status in separate neighbourhoods (mediated by planning), and ethnicity, where people of similar ethnic backgrounds try to co-locate.
About the same time, the Australian Bureau of Statistics defined distinct indexes of advantage and disadvantage taking the first principal component of sets of key variables that were thought to be important. These SEIFA indexes are regularly published for various jurisdictions, and are used frequently in spatial analysis.
Development indexes
PCA can be used as a formal method for the development of indexes. As an alternative confirmatory composite analysis has been proposed to develop and assess indexes.
The City Development Index was developed by PCA from about 200 indicators of city outcomes in a 1996 survey of 254 global cities. The first principal component was subject to iterative regression, adding the original variables singly until about 90% of its variation was accounted for. The index ultimately used about 15 indicators but was a good predictor of many more variables. Its comparative value agreed very well with a subjective assessment of the condition of each city. The coefficients on items of infrastructure were roughly proportional to the average costs of providing the underlying services, suggesting the Index was actually a measure of effective physical and social investment in the city.
The country-level Human Development Index (HDI) from UNDP, which has been published since 1990 and is very extensively used in development studies, has very similar coefficients on similar indicators, strongly suggesting it was originally constructed using PCA.
Population genetics
In 1978 Cavalli-Sforza and others pioneered the use of principal components analysis (PCA) to summarise data on variation in human gene frequencies across regions. The components showed distinctive patterns, including gradients and sinusoidal waves. They interpreted these patterns as resulting from specific ancient migration events.
Since then, PCA has been ubiquitous in population genetics, with thousands of papers using PCA as a display mechanism. Genetics varies largely according to proximity, so the first two principal components actually show spatial distribution and may be used to map the relative geographical location of different population groups, thereby showing individuals who have wandered from their original locations.
PCA in genetics has been technically controversial, in that the technique has been performed on discrete non-normal variables and often on binary allele markers. The lack of any measures of standard error in PCA are also an impediment to more consistent usage. In August 2022, the molecular biologist Eran Elhaik published a theoretical paper in Scientific Reports analyzing 12 PCA applications. He concluded that it was easy to manipulate the method, which, in his view, generated results that were 'erroneous, contradictory, and absurd.' Specifically, he argued, the results achieved in population genetics were characterized by cherry-picking and circular reasoning.
Market research and indexes of attitude
Market research has been an extensive user of PCA. It is used to develop customer satisfaction or customer loyalty scores for products, and with clustering, to develop market segments that may be targeted with advertising campaigns, in much the same way as factorial ecology will locate geographical areas with similar characteristics.
PCA rapidly transforms large amounts of data into smaller, easier-to-digest variables that can be more rapidly and readily analyzed. In any consumer questionnaire, there are series of questions designed to elicit consumer attitudes, and principal components seek out latent variables underlying these attitudes. For example, the Oxford Internet Survey in 2013 asked 2000 people about their attitudes and beliefs, and from these analysts extracted four principal component dimensions, which they identified as 'escape', 'social networking', 'efficiency', and 'problem creating'.
Another example from Joe Flood in 2008 extracted an attitudinal index toward housing from 28 attitude questions in a national survey of 2697 households in Australia. The first principal component represented a general attitude toward property and home ownership. The index, or the attitude questions it embodied, could be fed into a General Linear Model of tenure choice. The strongest determinant of private renting by far was the attitude index, rather than income, marital status or household type.
Quantitative finance
In quantitative finance, PCA is used
in financial risk management, and has been applied to other problems such as portfolio optimization.
PCA is commonly used in problems involving fixed income securities and portfolios, and interest rate derivatives.
Valuations here depend on the entire yield curve, comprising numerous highly correlated instruments, and PCA is used to define a set of components or factors that explain rate movements,
thereby facilitating the modelling.
One common risk management application is to calculating value at risk, VaR, applying PCA to the Monte Carlo simulation.
Here, for each simulation-sample, the components are stressed, and rates, and in turn option values, are then reconstructed;
with VaR calculated, finally, over the entire run.
PCA is also used in hedging exposure to interest rate risk, given partial durations and other sensitivities.
Under both, the first three, typically, principal components of the system are of interest (representing "shift", "twist", and "curvature").
These principal components are derived from an eigen-decomposition of the covariance matrix of yield at predefined maturities;
and where the variance of each component is its eigenvalue (and as the components are orthogonal, no correlation need be incorporated in subsequent modelling).
For equity, an optimal portfolio is one where the expected return is maximized for a given level of risk, or alternatively, where risk is minimized for a given return; see Markowitz model for discussion.
Thus, one approach is to reduce portfolio risk, where allocation strategies are applied to the "principal portfolios" instead of the underlying stocks.
A second approach is to enhance portfolio return, using the principal components to select companies' stocks with upside potential.
PCA has also been used to understand relationships between international equity markets, and within markets between groups of companies in industries or sectors.
PCA may also be applied to stress testing, essentially an analysis of a bank's ability to endure a hypothetical adverse economic scenario. Its utility is in "distilling the information contained in [several] macroeconomic variables into a more manageable data set, which can then [be used] for analysis." Here, the resulting factors are linked to e.g. interest rates – based on the largest elements of the factor's eigenvector – and it is then observed how a "shock" to each of the factors affects the implied assets of each of the banks.
Neuroscience
A variant of principal components analysis is used in neuroscience to identify the specific properties of a stimulus that increases a neuron's probability of generating an action potential. This technique is known as spike-triggered covariance analysis. In a typical application an experimenter presents a white noise process as a stimulus (usually either as a sensory input to a test subject, or as a current injected directly into the neuron) and records a train of action potentials, or spikes, produced by the neuron as a result. Presumably, certain features of the stimulus make the neuron more likely to spike. In order to extract these features, the experimenter calculates the covariance matrix of the spike-triggered ensemble, the set of all stimuli (defined and discretized over a finite time window, typically on the order of 100 ms) that immediately preceded a spike. The eigenvectors of the difference between the spike-triggered covariance matrix and the covariance matrix of the prior stimulus ensemble (the set of all stimuli, defined over the same length time window) then indicate the directions in the space of stimuli along which the variance of the spike-triggered ensemble differed the most from that of the prior stimulus ensemble. Specifically, the eigenvectors with the largest positive eigenvalues correspond to the directions along which the variance of the spike-triggered ensemble showed the largest positive change compared to the variance of the prior. Since these were the directions in which varying the stimulus led to a spike, they are often good approximations of the sought after relevant stimulus features.
In neuroscience, PCA is also used to discern the identity of a neuron from the shape of its action potential. Spike sorting is an important procedure because extracellular recording techniques often pick up signals from more than one neuron. In spike sorting, one first uses PCA to reduce the dimensionality of the space of action potential waveforms, and then performs clustering analysis to associate specific action potentials with individual neurons.
PCA as a dimension reduction technique is particularly suited to detect coordinated activities of large neuronal ensembles. It has been used in determining collective variables, that is, order parameters, during phase transitions in the brain.
Relation with other methods
Correspondence analysis
Correspondence analysis (CA)
was developed by Jean-Paul Benzécri
and is conceptually similar to PCA, but scales the data (which should be non-negative) so that rows and columns are treated equivalently. It is traditionally applied to contingency tables.
CA decomposes the chi-squared statistic associated to this table into orthogonal factors.
Because CA is a descriptive technique, it can be applied to tables for which the chi-squared statistic is appropriate or not.
Several variants of CA are available including detrended correspondence analysis and canonical correspondence analysis. One special extension is multiple correspondence analysis, which may be seen as the counterpart of principal component analysis for categorical data.
Factor analysis
Principal component analysis creates variables that are linear combinations of the original variables. The new variables have the property that the variables are all orthogonal. The PCA transformation can be helpful as a pre-processing step before clustering. PCA is a variance-focused approach seeking to reproduce the total variable variance, in which components reflect both common and unique variance of the variable. PCA is generally preferred for purposes of data reduction (that is, translating variable space into optimal factor space) but not when the goal is to detect the latent construct or factors.
Factor analysis is similar to principal component analysis, in that factor analysis also involves linear combinations of variables. Different from PCA, factor analysis is a correlation-focused approach seeking to reproduce the inter-correlations among variables, in which the factors "represent the common variance of variables, excluding unique variance". In terms of the correlation matrix, this corresponds with focusing on explaining the off-diagonal terms (that is, shared co-variance), while PCA focuses on explaining the terms that sit on the diagonal. However, as a side result, when trying to reproduce the on-diagonal terms, PCA also tends to fit relatively well the off-diagonal correlations. Results given by PCA and factor analysis are very similar in most situations, but this is not always the case, and there are some problems where the results are significantly different. Factor analysis is generally used when the research purpose is detecting data structure (that is, latent constructs or factors) or causal modeling. If the factor model is incorrectly formulated or the assumptions are not met, then factor analysis will give erroneous results.
-means clustering
It has been asserted that the relaxed solution of -means clustering, specified by the cluster indicators, is given by the principal components, and the PCA subspace spanned by the principal directions is identical to the cluster centroid subspace. However, that PCA is a useful relaxation of -means clustering was not a new result, and it is straightforward to uncover counterexamples to the statement that the cluster centroid subspace is spanned by the principal directions.
Non-negative matrix factorization
Non-negative matrix factorization (NMF) is a dimension reduction method where only non-negative elements in the matrices are used, which is therefore a promising method in astronomy, in the sense that astrophysical signals are non-negative. The PCA components are orthogonal to each other, while the NMF components are all non-negative and therefore constructs a non-orthogonal basis.
In PCA, the contribution of each component is ranked based on the magnitude of its corresponding eigenvalue, which is equivalent to the fractional residual variance (FRV) in analyzing empirical data. For NMF, its components are ranked based only on the empirical FRV curves. The residual fractional eigenvalue plots, that is, as a function of component number given a total of components, for PCA have a flat plateau, where no data is captured to remove the quasi-static noise, then the curves drop quickly as an indication of over-fitting (random noise). The FRV curves for NMF is decreasing continuously when the NMF components are constructed sequentially, indicating the continuous capturing of quasi-static noise; then converge to higher levels than PCA, indicating the less over-fitting property of NMF.
Iconography of correlations
It is often difficult to interpret the principal components when the data include many variables of various origins, or when some variables are qualitative. This leads the PCA user to a delicate elimination of several variables. If observations or variables have an excessive impact on the direction of the axes, they should be removed and then projected as supplementary elements. In addition, it is necessary to avoid interpreting the proximities between the points close to the center of the factorial plane.
The iconography of correlations, on the contrary, which is not a projection on a system of axes, does not have these drawbacks. We can therefore keep all the variables.
The principle of the diagram is to underline the "remarkable" correlations of the correlation matrix, by a solid line (positive correlation) or dotted line (negative correlation).
A strong correlation is not "remarkable" if it is not direct, but caused by the effect of a third variable. Conversely, weak correlations can be "remarkable". For example, if a variable Y depends on several independent variables, the correlations of Y with each of them are weak and yet "remarkable".
Generalizations
Sparse PCA
A particular disadvantage of PCA is that the principal components are usually linear combinations of all input variables. Sparse PCA overcomes this disadvantage by finding linear combinations that contain just a few input variables. It extends the classic method of principal component analysis (PCA) for the reduction of dimensionality of data by adding sparsity constraint on the input variables.
Several approaches have been proposed, including
a regression framework,
a convex relaxation/semidefinite programming framework,
a generalized power method framework
an alternating maximization framework
forward-backward greedy search and exact methods using branch-and-bound techniques,
Bayesian formulation framework.
The methodological and theoretical developments of Sparse PCA as well as its applications in scientific studies were recently reviewed in a survey paper.
Nonlinear PCA
Most of the modern methods for nonlinear dimensionality reduction find their theoretical and algorithmic roots in PCA or K-means. Pearson's original idea was to take a straight line (or plane) which will be "the best fit" to a set of data points. Trevor Hastie expanded on this concept by proposing Principal curves as the natural extension for the geometric interpretation of PCA, which explicitly constructs a manifold for data approximation followed by projecting the points onto it. | Mathematics | Statistics | null |
76347 | https://en.wikipedia.org/wiki/Hernia | Hernia | A hernia (: hernias or herniae, from Latin, meaning 'rupture') is the abnormal exit of tissue or an organ, such as the bowel, through the wall of the cavity in which it normally resides. The term is also used for the normal development of the intestinal tract, referring to the retraction of the intestine from the extra-embryonal navel coelom into the abdomen in the healthy embryo at about 7 weeks.
Various types of hernias can occur, most commonly involving the abdomen, and specifically the groin. Groin hernias are most commonly inguinal hernias but may also be femoral hernias. Other types of hernias include hiatus, incisional, and umbilical hernias. Symptoms are present in about 66% of people with groin hernias. This may include pain or discomfort in the lower abdomen, especially with coughing, exercise, or urinating or defecating. Often, it gets worse throughout the day and improves when lying down. A bulge may appear at the site of hernia, that becomes larger when bending down.
Groin hernias occur more often on the right than left side. The main concern is bowel strangulation, where the blood supply to part of the bowel is blocked. This usually produces severe pain and tenderness in the area. Hiatus, or hiatal hernias often result in heartburn but may also cause chest pain or pain while eating.
Risk factors for the development of a hernia include smoking, chronic obstructive pulmonary disease, obesity, pregnancy, peritoneal dialysis, collagen vascular disease and previous open appendectomy, among others. Predisposition to hernias is genetic and occur more often in certain families. Deleterious mutations causing predisposition to hernias seem to have dominant inheritance (especially for men). It is unclear if groin hernias are associated with heavy lifting. Hernias can often be diagnosed based on signs and symptoms. Occasionally, medical imaging is used to confirm the diagnosis or rule out other possible causes. The diagnosis of hiatus hernias is often done by endoscopy.
Groin hernias that do not cause symptoms in males do not need immediate surgical repair, a practice referred to as "watchful waiting". However most men tend to eventually undergo groin hernia surgery due to the development of pain. For women, however, repair is generally recommended due to the higher rate of femoral hernias, which have more complications. If strangulation occurs, immediate surgery is required. Repair may be done by open surgery, laparoscopic surgery, or robotic-assisted surgery. Open surgery has the benefit of possibly being done under local anesthesia rather than general anesthesia. Laparoscopic surgery generally has less pain following the procedure. A hiatus hernia may be treated with lifestyle changes such as raising the head of the bed, weight loss and adjusting eating habits. The medications H2 blockers or proton pump inhibitors may help. If the symptoms do not improve with medications, a surgery known as laparoscopic Nissen fundoplication may be an option.
Globally in 2019, there were 32.53 million prevalent cases of inguinal, femoral, and abdominal hernias, with a 95% uncertainty interval ranging from 27.71 to 37.79 million. Additionally, there were 13.02 million incident cases, with an uncertainty interval of 10.68 to 15.49 million. These figures reflect a 36.00% increase in prevalent cases and a 63.67% increase in incident cases compared to the numbers reported in 1990. About 27% of males and 3% of females develop a groin hernia at some point in their lives. Inguinal, femoral and abdominal hernias were present in 18.5 million people and resulted in 59,800 deaths in 2015. Groin hernias occur most often before the age of 1 and after the age of 50. It is not known how commonly hiatus hernias occur, with estimates in North America varying from 10% to 80%. The first known description of a hernia dates back to at least 1550 BC, in the Ebers Papyrus from Egypt.
Pathogenesis
Most hernias happen when the muscles and tendons in the belly weaken or get damaged, which makes it hard for them to keep the insides in place and support the body properly. The belly and pelvis act like a container made of muscles, tendons and bones. When pressure builds up inside this container, the muscles push back to keep everything in place. If the pressure gets too high, it may cause the belly's wall to break, leading to a hernia. Once a hernia starts, it keeps enlarging, because the tension on the wall there increases.
Epidemiology
About 27% of males and 3% of females develop a groin hernia at some time in their lives. In 2013 about 25 million people had a hernia. Inguinal, femoral and abdominal hernias resulted in 32,500 deaths globally in 2013 and 50,500 in 1990. Healthcare costs associated with abdominal wall hernias account for an annual expenditure of approximately 2.5 to 3 billion dollars.
Signs and symptoms
Symptoms and signs vary depending on the type of hernia. By far the most common hernias develop in the abdomen when a weakness in the abdominal wall evolves into a localized hole, or "defect", through which adipose tissue, or abdominal organs covered with peritoneum, may protrude. Another common hernia involves the spinal discs and causes sciatica. A hiatus hernia occurs when the stomach protrudes into the mediastinum through the esophageal opening in the diaphragm.
Hernias might manifest with pain in the area, a noticeable lump, or less specific symptoms caused by pressure on an organ stuck within the hernia, potentially leading to organ dysfunction. Typically, fatty tissue is the initial entrant into a hernia, but it might also involve an organ. Hernias are caused by a disruption or opening in the fascia, or fibrous tissue, which forms the abdominal wall. It is possible for the bulge associated with a hernia to come and go, but the defect in the tissue will persist.
Symptoms may or may not be present in some inguinal hernias. In the case of reducible hernias, a bulge in the groin or in another abdominal area can often be seen and felt. When standing, such a bulge becomes more obvious. Besides the bulge, other symptoms include pain in the groin that may also include a heavy or dragging sensation, and in men, there is sometimes pain and swelling in the scrotum around the testicular area.
Irreducible abdominal hernias or incarcerated hernias may be painful, but their most relevant symptom is that they cannot return to the abdominal cavity when pushed in. They may be chronic, although painless, and can lead to strangulation (loss of blood supply), obstruction (kinking of intestine), or both. Strangulated hernias are always painful and pain is followed by tenderness. Nausea, vomiting, or fever may occur in these cases due to bowel obstruction. Also, the hernia bulge, in this case, may turn red, purple or dark and pink.
In the diagnosis of abdominal hernias, imaging is the principal means of detecting internal diaphragmatic and other nonpalpable or unsuspected hernias. Multidetector CT (MDCT) can show with precision the anatomic site of the hernia sac, the contents of the sac, and any complications. MDCT also offers clear detail of the abdominal wall allowing wall hernias to be identified accurately.
Complications
Untreated hernia may be complicated by:
Inflammation
Obstruction of any lumen, such as bowel obstruction in intestinal hernias
Strangulation
Hydrocele of the hernial sac
Hemorrhage
Autoimmune problems
Irreducibility or incarceration, in which it cannot be reduced, or pushed back into place, at least not without very much external effort. In intestinal hernias, this also substantially increases the risk of bowel obstruction and strangulation.
Causes
Causes of hiatus hernia vary depending on each individual. Among the multiple causes, however, are the mechanical causes which include: improper heavy weight lifting, hard coughing bouts, sharp blows to the abdomen, and incorrect posture.
Furthermore, conditions that increase the pressure of the abdominal cavity may also cause hernias or worsen the existing ones. Some examples would be: obesity, straining during a bowel movement or urination (constipation, enlarged prostate), chronic lung disease, and also, fluid in the abdominal cavity (ascites).
Also, if muscles are weakened due to poor nutrition, smoking, and overexertion, hernias are more likely to occur.
The physiological school of thought contends that in the case of inguinal hernia, the above-mentioned are only an anatomical symptom of the underlying physiological cause. They contend that the risk of hernia is due to a physiological difference between patients who have hernia and those who do not, namely the presence of aponeurotic extensions from the transversus abdominis aponeurotic arch. There isn't any proof that being physically active will cause a hernia to get stuck or make an existing hernia worse.
Abdominal wall hernia may occur due to trauma. If this type of hernia is due to blunt trauma it is an emergency condition and could be associated with various solid organs and hollow viscus injuries.
Diagnosis
Inguinal
By far the most common hernias (up to 75% of all abdominal hernias) are inguinal hernias, which are further divided into the more common indirect inguinal hernia (2/3, depicted here), in which the inguinal canal is entered via a congenital weakness at its entrance (the internal inguinal ring), and the direct inguinal hernia type (1/3), where the hernia contents push through a weak spot in the back wall of the inguinal canal. An indirect inguinal hernia and a direct inguinal hernia can be distinguished by their positioning in relation to the inferior epigastric vessels. An indirect hernia is situated laterally to these vessels, whereas a direct hernia is positioned medially to them. Inguinal hernias are the most common type of hernia in both men and women. In some selected cases, they may require surgery.
There are special cases where a direct and indirect hernia appear together. A pantaloon hernia (or saddlebag hernia) is a combined direct and indirect hernia when the hernial sac protrudes on either side of the inferior epigastric vessels.
Additionally, though very rare, two or more indirect hernias may appear together such as in a double indirect hernia.
Femoral
Femoral hernias occur just below the inguinal ligament, when abdominal contents pass into the weak area at the posterior wall of the femoral canal. They can be hard to distinguish from the inguinal type (especially when ascending cephalad): however, they generally appear more rounded, and, in contrast to inguinal hernias, there is a strong female preponderance in femoral hernias. The incidence of strangulation in femoral hernias is high. Repair techniques are similar for femoral and inguinal hernia.
A Cooper's hernia is a femoral hernia with two sacs, the first being in the femoral canal, and the second passing through a defect in the superficial fascia and appearing almost immediately beneath the skin.
Umbilical
They involve protrusion of intra-abdominal contents through a weakness at the site of passage of the umbilical cord through the abdominal wall.
Umbilical hernias in adults are largely acquired, and are more frequent in obese or pregnant women. Abnormal decussation of fibers at the linea alba may be a contributing factor.
Incisional
An incisional hernia occurs when the defect is the result of an incompletely healed surgical wound. When these occur in median laparotomy incisions in the linea alba, they are termed ventral hernias. These occur in about 13% of people at 2 years following surgery.
Diaphragmatic
Higher in the abdomen, an (internal) "diaphragmatic hernia" results when part of the stomach or intestine protrudes into the chest cavity through a defect in the diaphragm.
A hiatus hernia is a particular variant of this type, in which the normal passageway through which the esophagus meets the stomach (esophageal hiatus) serves as a functional "defect", allowing part of the stomach to (periodically) "herniate" into the chest. Hiatus hernias may be either "sliding", in which the gastroesophageal junction itself slides through the defect into the chest, or non-sliding (also known as para-esophageal), in which case the junction remains fixed while another portion of the stomach moves up through the defect. Non-sliding or para-esophageal hernias can be dangerous as they may allow the stomach to rotate and obstruct. Repair is usually advised.
A congenital diaphragmatic hernia is a distinct problem, occurring in up to 1 in 2000 births, and requiring pediatric surgery. Intestinal organs may herniate through several parts of the diaphragm, posterolateral (in Bochdalek's triangle (lumbocostal triangle), resulting in a Bochdalek hernia), or anteromedial-retrosternal (in the cleft of foramina of Morgagni (sternocostal triangle), resulting in a Morgagni's hernia).
Other hernias
Since many organs or parts of organs can herniate through many orifices, it is very difficult to give an exhaustive list of hernias, with all synonyms and eponyms. The above article deals mostly with "visceral hernias", where the herniating tissue arises within the abdominal cavity. Other hernia types and unusual types of visceral hernias are listed below, in alphabetical order:
Abdominal wall hernias:
Umbilical hernia
Epigastric hernia: a hernia through the linea alba above the umbilicus.
Spigelian hernia, also known as spontaneous lateral ventral hernia
Amyand's hernia: containing the appendix vermiformis within the hernia sac
Brain herniation, sometimes referred to as brain hernia, is a potentially deadly side effect of very high intracranial pressure that occurs when a part of the brain is squeezed across structures within the skull.
Broad ligament hernia, of the uterus.
Double indirect hernia: an indirect inguinal hernia with two hernia sacs, without a concomitant direct hernia component (as seen in a pantaloon hernia).
Hiatus hernia: a hernia due to "short oesophagus" — insufficient elongation — stomach is displaced into the thorax
Littre's hernia: a hernia involving a Meckel's diverticulum. It is named after the French anatomist Alexis Littré (1658–1726).
Lumbar hernia: a hernia in the lumbar region (not to be confused with a lumbar disc hernia), contains the following entities:
Petit's hernia: a hernia through Petit's triangle (inferior lumbar triangle). It is named after French surgeon Jean Louis Petit (1674–1750).
Grynfeltt's hernia: a hernia through Grynfeltt-Lesshaft triangle (superior lumbar triangle). It is named after physician Joseph Grynfeltt (1840–1913).
Maydl's hernia: two adjacent loops of small intestine are within a hernial sac with a tight neck. The intervening portion of bowel within the abdomen is deprived of its blood supply and eventually becomes necrotic.
Obturator hernia: hernia through obturator canal
Parastomal hernias, which is when tissue protrudes adjacent to a stoma tract.
Paraumbilical hernia: a type of umbilical hernia occurring in adults
Perineal hernia: a perineal hernia protrudes through the muscles and fascia of the perineal floor. It may be primary but usually is acquired following perineal prostatectomy, abdominoperineal resection of the rectum, or pelvic exenteration.
Properitoneal hernia: rare hernia located directly above the peritoneum, for example, when part of inguinal hernia projects from the deep inguinal ring to the preperitoneal space.
Retrocolic hernia: entrapment of portions of the small intestine behind the mesocolon.
Richter's hernia: a hernia involving only one sidewall of the bowel, which can result in bowel strangulation leading to perforation through ischaemia without causing bowel obstruction or any of its warning signs. It is named after German surgeon August Gottlieb Richter (1742–1812).
Sliding hernia: occurs when an organ drags along part of the peritoneum, or, in other words, the organ is part of the hernia sac. The colon and the urinary bladder are often involved. The term also frequently refers to sliding hernias of the stomach.
Sciatic hernia: this hernia in the greater sciatic foramen most commonly presents as an uncomfortable mass in the gluteal area. Bowel obstruction may also occur. This type of hernia is only a rare cause of sciatic neuralgia.
Sports hernia: a hernia characterized by chronic groin pain in athletes and a dilated superficial inguinal ring.
Tibialis anterior hernia: can present as a bulge in the shins. Pain on rest, walking, or during exercise may occur. The bulge can typically not be present unless pressure or flexing of the leg occurs.
Velpeau hernia: a hernia in the groin in front of the femoral blood vessels
Treatment
Truss
The benefits of the use of an external device to maintain reduction of the hernia without repairing the underlying defect (such as hernia trusses, trunks, belts, etc.) are unclear.
Surgery
Surgery is recommended for some types of hernias to prevent complications such as obstruction of the bowel or strangulation of the tissue, although umbilical hernias and hiatus hernias may be watched, or are treated with medication. Most abdominal hernias can be surgically repaired, but surgery has complications. Prior to surgery patients should be medically optimized receive guidance about changing factors that can be controlled, such as quitting smoking, managing medical conditions like diabetes effectively, and working on losing weight.
Three primary methods can be utilized: open surgery, laparoscopy, or robotic techniques. Fixing an inguinal hernia using laparoscopy causes less pain, speeds up recovery, and shows similar low rates of the hernia coming back compared to the traditional open repair method. However, open surgery can be done sometimes without general anesthesia. Using local anesthesia for open groin hernia repair, particularly in patients with additional health issues, leads to fewer complications and reduced costs. Studies show that compared to regional or general anesthesia, local anesthesia results in less postoperative pain, shorter recovery times, and decreased unplanned overnight stays. However, it might not be enough for repairing large hernias or in patients with abdominal domain loss, where general anesthesia is preferred.
Robot-assisted hernia surgery has also recently gained popularity as safe alternatives to open surgery. Robotic surgery for inguinal hernia repair shows outcomes comparable to laparoscopic surgery. The rates of overall complications, long-lasting postoperative pain, urinary retention, and 30-day re-admission are very similar between these two methods. Just like in other areas of general surgery, it has been noted that robotic surgery for inguinal hernia repair takes more time in the operating room compared to the laparoscopic approach.
Uncomplicated hernias are principally repaired by pushing back, or "reducing", the herniated tissue, and then mending the weakness in muscle tissue (an operation called herniorrhaphy). If complications have occurred, the surgeon will check the viability of the herniated organ and remove part of it if necessary.
Muscle reinforcement techniques often involve synthetic materials (a mesh prosthesis). The mesh is placed either over the defect (anterior repair) or under the defect (posterior repair). At times staples are used to keep the mesh in place. These mesh repair methods are often called "tension free" repairs because, unlike some suture methods (e.g., Shouldice), muscle is not pulled together under tension. However, this widely used terminology is misleading, as there are many tension-free suture methods that do not use mesh (e.g., Desarda, Guarnieri, Lipton-Estrin, etc.).
Evidence suggests that tension-free methods (with or without mesh) often have lower percentage of recurrences and the fastest recovery period compared to tension suture methods. However, the use of prosthetic mesh appears to have a higher likelihood of causing long-term pain and can also lead to infections.
The frequency of surgical correction ranges from 10 per 100,000 (U.K.) to 28 per 100,000 (U.S.). After elective surgery, the 30-day mortality rate for inguinal or femoral hernia repair stands at 0.1 percent, but it increases to 2.8 to 3.1 percent after urgent surgery. When a bowel resection is part of the hernia repair, the mortality rate is even higher. Older age, femoral hernias, female sex, and urgent repair are identified as other factors linked to a higher risk of mortality.
Post-Operative Complications
Some complications from surgery in order of prevalence include a seroma/hematoma formation, urinary retention, neuralgias, testicular pain/swelling, mesh infection/wound infection, and recurrence. A seroma is often seen after an indirect hernia repair and resolves spontaneously over 4–6 weeks. To prevent a seroma it's important to reduce the amount of cutting around the hernia sac where it's connected to the cord structures. Additionally, securely attaching the hernia sac to the pubic bone and creating small openings in the tissue around a direct hernia can help. In cases of heavy bleeding or extensive cutting, certain surgeons may opt to insert a drain. Urinary retention is often seen in elderly patients, these patients can be catheterized prior to surgery if there is a risk. Other complications may arise post-operatively, including rejection of the mesh that is used to repair the hernia. In the event of a mesh rejection, the mesh will very likely need to be removed. Mesh rejection can be detected by obvious, sometimes localized swelling and pain around the mesh area. Continuous discharge from the scar is likely for a while after the mesh has been removed. A surgically treated hernia can lead to complications such as inguinodynia.
Recovery
Many patients are managed through day surgery centers and are able to return to work within a week or two, though intense activities are prohibited for a longer period. People who have their hernias repaired with mesh often recover within a month, but pain can last longer. Surgical complications may include pain that lasts more than three months, surgical site infections, nerve and blood vessel injuries, injury to nearby organs, and hernia recurrence. Pain that lasts more than three months occurs in about 10% of people following hernia repair.
External links
Disorders of fascia
Congenital disorders of musculoskeletal system
Wikipedia medicine articles ready to translate
Acute pain
Wikipedia emergency medicine articles ready to translate | Biology and health sciences | Types | Health |
76408 | https://en.wikipedia.org/wiki/Transverse%20wave | Transverse wave | In physics, a transverse wave is a wave that oscillates perpendicularly to the direction of the wave's advance. In contrast, a longitudinal wave travels in the direction of its oscillations. All waves move energy from place to place without transporting the matter in the transmission medium if there is one. Electromagnetic waves are transverse without requiring a medium. The designation “transverse” indicates the direction of the wave is perpendicular to the displacement of the particles of the medium through which it passes, or in the case of EM waves, the oscillation is perpendicular to the direction of the wave.
A simple example is given by the waves that can be created on a horizontal length of string by anchoring one end and moving the other end up and down. Another example is the waves that are created on the membrane of a drum. The waves propagate in directions that are parallel to the membrane plane, but each point in the membrane itself gets displaced up and down, perpendicular to that plane. Light is another example of a transverse wave, where the oscillations are the electric and magnetic fields, which point at right angles to the ideal light rays that describe the direction of propagation.
Transverse waves commonly occur in elastic solids due to the shear stress generated; the oscillations in this case are the displacement of the solid particles away from their relaxed position, in directions perpendicular to the propagation of the wave. These displacements correspond to a local shear deformation of the material. Hence a transverse wave of this nature is called a shear wave. Since fluids cannot resist shear forces while at rest, propagation of transverse waves inside the bulk of fluids is not possible. In seismology, shear waves are also called secondary waves or S-waves.
Transverse waves are contrasted with longitudinal waves, where the oscillations occur in the direction of the wave. The standard example of a longitudinal wave is a sound wave or "pressure wave" in gases, liquids, or solids, whose oscillations cause compression and expansion of the material through which the wave is propagating. Pressure waves are called "primary waves", or "P-waves" in geophysics.
Water waves involve both longitudinal and transverse motions.
Mathematical formulation
Mathematically, the simplest kind of transverse wave is a plane linearly polarized sinusoidal one. "Plane" here means that the direction of propagation is unchanging and the same over the whole medium; "linearly polarized" means that the direction of displacement too is unchanging and the same over the whole medium; and the magnitude of the displacement is a sinusoidal function only of time and of position along the direction of propagation.
The motion of such a wave can be expressed mathematically as follows. Let be the direction of propagation (a vector with unit length), and any reference point in the medium. Let be the direction of the oscillations (another unit-length vector perpendicular to d). The displacement of a particle at any point of the medium and any time t (seconds) will be
where A is the wave's amplitude or strength, T is its period, v is the speed of propagation, and is its phase at t = 0 seconds at . All these parameters are real numbers. The symbol "•" denotes the inner product of two vectors.
By this equation, the wave travels in the direction and the oscillations occur back and forth along the direction . The wave is said to be linearly polarized in the direction .
An observer that looks at a fixed point will see the particle there move in a simple harmonic (sinusoidal) motion with period T seconds, with maximum particle displacement A in each sense; that is, with a frequency of f = 1/T full oscillation cycles every second. A snapshot of all particles at a fixed time t will show the same displacement for all particles on each plane perpendicular to , with the displacements in successive planes forming a sinusoidal pattern, with each full cycle extending along by the wavelength λ = v T = v/f. The whole pattern moves in the direction with speed V.
The same equation describes a plane linearly polarized sinusoidal light wave, except that the "displacement" S(, t) is the electric field at point and time t. (The magnetic field will be described by the same equation, but with a "displacement" direction that is perpendicular to both and , and a different amplitude.)
Superposition principle
In a homogeneous linear medium, complex oscillations (vibrations in a material or light flows) can be described as the superposition of many simple sinusoidal waves, either transverse or longitudinal.
The vibrations of a violin string create standing waves, for example, which can be analyzed as the sum of many transverse waves of different frequencies moving in opposite directions to each other, that displace the string either up or down or left to right. The antinodes of the waves align in a superposition .
Circular polarization
If the medium is linear and allows multiple independent displacement directions for the same travel direction , we can choose two mutually perpendicular directions of polarization, and express any wave linearly polarized in any other direction as a linear combination (mixing) of those two waves.
By combining two waves with same frequency, velocity, and direction of travel, but with different phases and independent displacement directions, one obtains a circularly or elliptically polarized wave. In such a wave the particles describe circular or elliptical trajectories, instead of moving back and forth.
It may help understanding to revisit the thought experiment with a taut string mentioned above. Notice that you can also launch waves on the string by moving your hand to the right and left instead of up and down. This is an important point. There are two independent (orthogonal) directions that the waves can move. (This is true for any two directions at right angles, up and down and right and left are chosen for clarity.) Any waves launched by moving your hand in a straight line are linearly polarized waves.
But now imagine moving your hand in a circle. Your motion will launch a spiral wave on the string. You are moving your hand simultaneously both up and down and side to side. The maxima of the side to side motion occur a quarter wavelength (or a quarter of a way around the circle, that is 90 degrees or π/2 radians) from the maxima of the up and down motion. At any point along the string, the displacement of the string will describe the same circle as your hand, but delayed by the propagation speed of the wave. Notice also that you can choose to move your hand in a clockwise circle or a counter-clockwise circle. These alternate circular motions produce right and left circularly polarized waves.
To the extent your circle is imperfect, a regular motion will describe an ellipse, and produce elliptically polarized waves. At the extreme of eccentricity your ellipse will become a straight line, producing linear polarization along the major axis of the ellipse. An elliptical motion can always be decomposed into two orthogonal linear motions of unequal amplitude and 90 degrees out of phase, with circular polarization being the special case where the two linear motions have the same amplitude.
Power in a transverse wave in string
(Let the linear mass density of the string be μ.)
The kinetic energy of a mass element in a transverse wave is given by:
In one wavelength, kinetic energy
Using Hooke's law the potential energy in mass element
And the potential energy for one wavelength
So, total energy in one wavelength
Therefore average power is
| Physical sciences | Waves | Physics |
76653 | https://en.wikipedia.org/wiki/Naval%20architecture | Naval architecture | Naval architecture, or naval engineering, is an engineering discipline incorporating elements of mechanical, electrical, electronic, software and safety engineering as applied to the engineering design process, shipbuilding, maintenance, and operation of marine vessels and structures. Naval architecture involves basic and applied research, design, development, design evaluation (classification) and calculations during all stages of the life of a marine vehicle. Preliminary design of the vessel, its detailed design, construction, trials, operation and maintenance, launching and dry-docking are the main activities involved. Ship design calculations are also required for ships being modified (by means of conversion, rebuilding, modernization, or repair). Naval architecture also involves formulation of safety regulations and damage-control rules and the approval and certification of ship designs to meet statutory and non-statutory requirements.
Main subjects
The word "vessel" includes every description of watercraft, mainly ships and boats, but also including non-displacement craft, WIG craft and seaplanes, used or capable of being used as a means of transportation on water. The principal elements of naval architecture are detailed in the following sections.
Hydrostatics
Hydrostatics concerns the conditions to which the vessel is subjected while at rest in water and to its ability to remain afloat. This involves computing buoyancy, displacement, and other hydrostatic properties such as trim (the measure of the longitudinal inclination of the vessel) and stability (the ability of a vessel to restore itself to an upright position after being inclined by wind, sea, or loading conditions).
Hydrodynamics
Hydrodynamics concerns the flow of water around the ship's hull, bow, and stern, and over bodies such as propeller blades or rudder, or through thruster tunnels.
Ship resistance and propulsion concern resistance towards motion in water primarily caused due to flow of water around the hull. Powering calculation is done based on this.
Propulsion is used to move the vessel through water using propellers, thrusters, water jets, sails etc. Engine types are mainly internal combustion. Some vessels are electrically powered using nuclear or solar energy.
Ship motions involves motions of the vessel in seaway and its responses in waves and wind.
Controllability (maneuvering) involves controlling and maintaining position and direction of the vessel.
Flotation and stability
While atop a liquid surface a floating body has 6 degrees of freedom in its movements, these are categorized in either translation or rotation.
Translation
Sway: transverse
Surge: fore and aft
Heave: vertical
Rotation
Yaw: about a vertical axis
Pitch or trim: about a transverse axis
Roll or heel: about a fore and aft axis
Longitudinal stability for longitudinal inclinations, the stability depends upon the distance between the center of gravity and the longitudinal meta-center. In other words, the basis in which the ship maintains its center of gravity is its distance set equally apart from both the aft and forward section of the ship.
While a body floats on a liquid surface it still encounters the force of gravity pushing down on it. In order to stay afloat and avoid sinking there is an opposed force acting against the body known as the hydrostatic pressures. The forces acting on the body must be of the same magnitude and same line of motion in order to maintain the body at equilibrium. This description of equilibrium is only present when a freely floating body is in still water, when other conditions are present the magnitude of which these forces shifts drastically creating the swaying motion of the body.
The buoyancy force is equal to the weight of the body, in other words, the mass of the body is equal to the mass of the water displaced by the body. This adds an upward force to the body by the amount of surface area times the area displaced in order to create an equilibrium between the surface of the body and the surface of the water.
The stability of a ship under most conditions is able to overcome any form or restriction or resistance encountered in rough seas; however, ships have undesirable roll characteristics when the balance of oscillations in roll is two times that of oscillations in heave, thus causing the ship to capsize.
Structures
Structures involves selection of material of construction, structural analysis of global and local strength of the vessel, vibration of the structural components and structural responses of the vessel during motions in seaway. Depending on type of ship, the structure and design will vary in what material to use as well as how much of it. Some ships are made from glass reinforced plastics but the vast majority are steel with possibly some aluminium in the superstructure.
The complete structure of the ship is designed with panels shaped in a rectangular form consisting of steel plating supported on four edges. Combined in a large surface area the Grillages create the hull of the ship, deck, and bulkheads while still providing mutual support of the frames. Though the structure of the ship is sturdy enough to hold itself together the main force it has to overcome is longitudinal bending creating a strain against its hull, its structure must be designed so that the material is disposed as much forward and aft as possible.
The principal longitudinal elements are the deck, shell plating, inner bottom all of which are in the form of grillages, and additional longitudinal stretching to these. The dimensions of the ship are in order to create enough spacing between the stiffeners in prevention of buckling. Warships have used a longitudinal system of stiffening that many modern commercial vessels have adopted. This system was widely used in early merchant ships such as the SS Great Eastern, but later shifted to transversely framed structure another concept in ship hull design that proved more practical. This system was later implemented on modern vessels such as tankers because of its popularity and was then named the Isherwood System.
The arrangement of the Isherwood system consists of stiffening decks both side and bottom by longitudinal members, they are separated enough so they have the same distance between them as the frames and beams. This system works by spacing out the transverse members that support the longitudinal by about 3 or 4 meters, with the wide spacing this causes the traverse strength needed by displacing the amount of force the bulkheads provide.
Arrangements
Arrangements involves concept design, layout and access, fire protection, allocation of spaces, ergonomics and capacity.
Construction
Construction depends on the material used. When steel or aluminium is used this involves welding of the plates and profiles after rolling, marking, cutting and bending as per the structural design drawings or models, followed by erection and launching. Other joining techniques are used for other materials like fibre reinforced plastic and glass-reinforced plastic. The process of construction is thought-out cautiously while considering all factors like safety, strength of structure, hydrodynamics, and ship arrangement. Each factor considered presents a new option for materials to consider as well as ship orientation. When the strength of the structure is considered the acts of ship collision are considered in the way that the ships structure is altered. Therefore, the properties of materials are considered carefully as applied material on the struck ship has elastic properties, the energy absorbed by the ship being struck is then deflected in the opposite direction, so both ships go through the process of rebounding to prevent further damage.
Science and craft
Traditionally, naval architecture has been more craft than science. The suitability of a vessel's shape was judged by looking at a half-model of a vessel or a prototype. Ungainly shapes or abrupt transitions were frowned on as being flawed. This included rigging, deck arrangements, and even fixtures. Subjective descriptors such as ungainly, full, and fine were used as a substitute for the more precise terms used today. A vessel was, and still is described as having a ‘fair’ shape. The term ‘fair’ is meant to denote not only a smooth transition from fore to aft but also a shape that was ‘right.’ Determining what is ‘right’ in a particular situation in the absence of definitive supporting analysis encompasses the art of naval architecture to this day.
Modern low-cost digital computers and dedicated software, combined with extensive research to correlate full-scale, towing tank and computational data, have enabled naval architects to more accurately predict the performance of a marine vehicle. These tools are used for static stability (intact and damaged), dynamic stability, resistance, powering, hull development, structural analysis, green water modelling, and slamming analysis. Data are regularly shared in international conferences sponsored by RINA, Society of Naval Architects and Marine Engineers (SNAME) and others. Computational Fluid Dynamics is being applied to predict the response of a floating body in a random sea.
The naval architect
Due to the complexity associated with operating in a marine environment, naval architecture is a co-operative effort between groups of technically skilled individuals who are specialists in particular fields, often coordinated by a lead naval architect. This inherent complexity also means that the analytical tools available are much less evolved than those for designing aircraft, cars and even spacecraft. This is due primarily to the paucity of data on the environment the marine vehicle is required to work in and the complexity of the interaction of waves and wind on a marine structure.
A naval architect is an engineer who is responsible for the design, classification, survey, construction, and/or repair of ships, boats, other marine vessels, and offshore structures, both commercial and military, including:
Merchant ships – oil tankers, gas tankers, cargo ships, bulk carriers, container ships
Passenger/vehicle ferries, cruise ships
Warships – frigates, destroyers, aircraft carriers, amphibious ships
Submarines and underwater vehicles
Icebreakers
High speed craft – hovercraft, multi-hull ships, hydrofoil craft
Workboats – barges, fishing boats, anchor handling tug supply vessels, platform supply vessels, tug boats, pilot vessels, rescue craft
Yachts, power boats, and other recreational watercraft
Offshore platforms and subsea developments
Some of these vessels are amongst the largest (such as supertankers), most complex (such as aircraft carriers), and highly valued movable structures produced by mankind. They are typically the most efficient method of transporting the world's raw materials and products. Modern engineering on this scale is essentially a team activity conducted by specialists in their respective fields and disciplines.
Naval architects integrate these activities. This demanding leadership role requires managerial qualities and the ability to bring together the often-conflicting demands of the various design constraints to produce a product which is fit for the purpose.
In addition to this leadership role, a naval architect also has a specialist function in ensuring that a safe, economic, environmentally sound and seaworthy design is produced. To undertake all these tasks, a naval architect must have an understanding of many branches of engineering and must be in the forefront of high technology areas. He or she must be able to effectively utilize the services provided by scientists, lawyers, accountants, and business people of many kinds.
Naval architects typically work for shipyards, ship owners, design firms and consultancies, equipment manufacturers, Classification societies, regulatory bodies (Admiralty law), navies, and governments. A small majority of Naval Architects also work in education, of which only 5 universities in the United States are accredited with Naval Architecture & Marine Engineering programs. The United States Naval Academy is home to one of the most knowledgeable professors of Naval Architecture; CAPT. Michael Bito, USN.
| Technology | Naval transport | null |
76812 | https://en.wikipedia.org/wiki/Dipper | Dipper | Dippers are members of the genus Cinclus in the bird family Cinclidae, so-called because of their bobbing or dipping movements. They are unique among passerines for their ability to dive and swim underwater.
Taxonomy
The genus Cinclus was introduced by the German naturalist Moritz Balthasar Borkhausen in 1797 with the white-throated dipper (Cinclus cinclus) as the type species. The name cinclus is from the Ancient Greek word kinklos that was used to describe small tail-wagging birds that resided near water.
Cinclus is the only genus in the family Cinclidae. The white-throated dipper and American dipper are also known in Britain and America, respectively, as the water ouzel (sometimes spelt "ousel") – ouzel originally meant the only distantly related but superficially similar Eurasian blackbird (Old English osle). Ouzel also survives as the name of a relative of the blackbird, the ring ouzel.
The genus contains five species:
White-throated dipper or European dipper, Cinclus cinclus
Brown dipper Cinclus pallasii
American dipper Cinclus mexicanus
White-capped dipper Cinclus leucocephalus
Rufous-throated dipper Cinclus schulzii
A 2002 molecular phylogenetic study of the dippers looked at the DNA sequences of two mitochondrial genes. It found that the Eurasian white-throated dipper and brown dipper are sister species as are the South American white-capped dipper and rufous-throated dipper. The study also showed that the dipper family, Cinclidae, is most closely related to the thrush family, Turdidae.
Description
Dippers are small, chunky, stout, short-tailed, short-winged, strong-legged birds. The different species are generally dark brown (sometimes nearly black), or brown and white in colour, apart from the rufous-throated dipper, which is brown with a reddish-brown throat patch. Sizes range from in length and in weight, with males larger than females. Their short wings give them a distinctive whirring flight. They have a characteristic bobbing motion when perched beside the water, giving them their name. While under water, they are covered by a thin, silvery film of air, due to small bubbles being trapped on the surface of the plumage.
Distribution and habitat
Dippers are found in suitable freshwater habitats in the highlands of the Americas, Europe and Asia. In Africa they are only found in the Atlas Mountains of Morocco. They inhabit the banks of fast-moving upland rivers with cold, clear waters, though, outside the breeding season, they may visit lake shores and sea coasts.
Adaptations
Unlike many water birds, dippers are generally similar in form to many terrestrial birds (for example, they do not have webbed feet), but they do have some morphological and physiological adaptations to their aquatic habits. They have evolved solid bones to reduce their buoyancy, and their wings are relatively short but strongly muscled, enabling them to be used as flippers underwater. The plumage is dense, with a large preen gland for waterproofing their feathers. Relatively long legs and sharp claws enable them to hold on to rocks in swift water. Their eyes have well-developed focus muscles that can change the curvature of the lens to enhance underwater vision. They have nasal flaps to prevent water entering their nostrils.
The high haemoglobin concentration in their blood gives them a capacity to store oxygen greater than that of other birds, allowing them to remain underwater for 30 seconds or more, whilst their basal metabolic rate is approximately one-third slower than typical terrestrial passerines of similar mass. One small population wintering at a hot spring in Suntar-Khayata Mountains of Siberia feeds underwater when air temperatures drop below .
Behaviour
Food
Dippers forage for small animal prey in and along the margins of fast-flowing freshwater streams and rivers. They perch on rocks and feed at the edge of the water, but they often also grip the rocks firmly and walk down them beneath the water until partly or wholly submerged. They then search underwater for prey between and beneath stones and debris; they can also swim with their wings. The two South American species swim and dive less often than the three northern ones. Their prey consists primarily of invertebrates such as the nymphs or larvae of mayflies, blackflies, stoneflies and caddisflies, as well as small fish and fish eggs. Molluscs and crustaceans are also consumed, especially in winter when insect larvae are less available.
Breeding
Linear breeding territories are established by pairs of dippers along suitable rivers, and maintained against incursion by other dippers. Within their territory the pair must have a good nest site and roost sites, but the main factor affecting the length of the territory is the availability of sufficient food to feed themselves and their broods. Consequently, the length of a territory may vary from about to over .
Dipper nests are usually large, round, domed structures made of moss, with an internal cup of grass and rootlets, and a side entrance hole. They are often built in confined spaces over, or close to, running water. The site may be on a ledge or bank, in a crevice or drainpipe, or beneath a bridge. Tree sites are rare.
The usual clutch-size of the three northern dipper species is four or five; those of the South American species is not well known, though some evidence suggests that of the rufous-throated dipper is two. The incubation period of 16 or 17 days is followed by the hatching of altricial young which are brooded by the female alone for the next 12 to 13 days. The nestlings are fed by both parents and the whole fledging period is about 20–24 days. Young dippers usually become independent of their parents within a couple of weeks of leaving the nest. Dippers may raise second broods if conditions allow. The maximum recorded age from ring-recovery data of a white-throated dipper is 10 years and 7 months for a bird ringed in Finland. The maximum age for an American dipper is 8 years and 1 month for a bird ringed and recovered in South Dakota.
Communication
Dippers' calls are loud and high-pitched, being similar to calls made by other birds on fast rivers; the call frequencies lying within a narrow range of 4.0–6.5 kHz, well above the torrent noise frequency of maximum 2 kHz. Dippers also communicate visually by their characteristic dipping or bobbing movements, as well as by blinking rapidly to expose the white feathers on their upper eyelids as a series of white flashes in courtship and threat displays.
Conservation
Dippers are completely dependent on fast-flowing rivers with clear water, accessible food and secure nest-sites. They may be threatened by anything that affects these needs such as water pollution, acidification and turbidity caused by erosion. River regulation through the creation of dams and reservoirs, as well as channelization, can degrade and destroy dipper habitat.
Dippers are also sometimes hunted or otherwise persecuted by humans for various reasons. The Cyprus race of the white-throated dipper is extinct. In the Atlas Mountains dippers are claimed to have aphrodisiacal properties. In parts of Scotland and Germany, until the beginning of the 20th century, bounties were paid for killing dippers because of a misguided perception that they were detrimental to fish stocks through predation on the eggs and fry of salmonids.
Despite threats to local populations, the conservation status of most dipper species is considered to be of least concern. The one exception, the rufous-throated dipper, is classified as vulnerable because of its small, fragmented and declining population which is threatened, especially in Argentina, by changes in river management.
| Biology and health sciences | Passerida | Animals |
76848 | https://en.wikipedia.org/wiki/Periodic%20function | Periodic function | A periodic function, also called a periodic waveform (or simply periodic wave), is a function that repeats its values at regular intervals or periods. The repeatable part of the function or waveform is called a cycle. For example, the trigonometric functions, which repeat at intervals of radians, are periodic functions. Periodic functions are used throughout science to describe oscillations, waves, and other phenomena that exhibit periodicity. Any function that is not periodic is called aperiodic.
Definition
A function is said to be periodic if, for some nonzero constant , it is the case that
for all values of in the domain. A nonzero constant for which this is the case is called a period of the function. If there exists a least positive constant with this property, it is called the fundamental period (also primitive period, basic period, or prime period.) Often, "the" period of a function is used to mean its fundamental period. A function with period will repeat on intervals of length , and these intervals are sometimes also referred to as periods of the function.
Geometrically, a periodic function can be defined as a function whose graph exhibits translational symmetry, i.e. a function is periodic with period if the graph of is invariant under translation in the -direction by a distance of . This definition of periodicity can be extended to other geometric shapes and patterns, as well as be generalized to higher dimensions, such as periodic tessellations of the plane. A sequence can also be viewed as a function defined on the natural numbers, and for a periodic sequence these notions are defined accordingly.
Examples
Real number examples
The sine function is periodic with period , since
for all values of . This function repeats on intervals of length (see the graph to the right).
Everyday examples are seen when the variable is time; for instance the hands of a clock or the phases of the moon show periodic behaviour. Periodic motion is motion in which the position(s) of the system are expressible as periodic functions, all with the same period.
For a function on the real numbers or on the integers, that means that the entire graph can be formed from copies of one particular portion, repeated at regular intervals.
A simple example of a periodic function is the function that gives the "fractional part" of its argument. Its period is 1. In particular,
The graph of the function is the sawtooth wave.
The trigonometric functions sine and cosine are common periodic functions, with period (see the figure on the right). The subject of Fourier series investigates the idea that an 'arbitrary' periodic function is a sum of trigonometric functions with matching periods.
According to the definition above, some exotic functions, for example the Dirichlet function, are also periodic; in the case of Dirichlet function, any nonzero rational number is a period.
Complex number examples
Using complex variables we have the common period function:
Since the cosine and sine functions are both periodic with period , the complex exponential is made up of cosine and sine waves. This means that Euler's formula (above) has the property such that if is the period of the function, then
Double-periodic functions
A function whose domain is the complex numbers can have two incommensurate periods without being constant. The elliptic functions are such functions. ("Incommensurate" in this context means not real multiples of each other.)
Properties
Periodic functions can take on values many times. More specifically, if a function is periodic with period , then for all in the domain of and all positive integers ,
If is a function with period , then , where is a non-zero real number such that is within the domain of , is periodic with period . For example, has period and, therefore, will have period .
Some periodic functions can be described by Fourier series. For instance, for L2 functions, Carleson's theorem states that they have a pointwise (Lebesgue) almost everywhere convergent Fourier series. Fourier series can only be used for periodic functions, or for functions on a bounded (compact) interval. If is a periodic function with period that can be described by a Fourier series, the coefficients of the series can be described by an integral over an interval of length .
Any function that consists only of periodic functions with the same period is also periodic (with period equal or smaller), including:
addition, subtraction, multiplication and division of periodic functions, and
taking a power or a root of a periodic function (provided it is defined for all ).
Generalizations
Antiperiodic functions
One subset of periodic functions is that of antiperiodic functions. This is a function such that for all . For example, the sine and cosine functions are -antiperiodic and -periodic. While a -antiperiodic function is a -periodic function, the converse is not necessarily true.
Bloch-periodic functions
A further generalization appears in the context of Bloch's theorems and Floquet theory, which govern the solution of various periodic differential equations. In this context, the solution (in one dimension) is typically a function of the form
where is a real or complex number (the Bloch wavevector or Floquet exponent). Functions of this form are sometimes called Bloch-periodic in this context. A periodic function is the special case , and an antiperiodic function is the special case . Whenever is rational, the function is also periodic.
Quotient spaces as domain
In signal processing you encounter the problem, that Fourier series represent periodic functions and that Fourier series satisfy convolution theorems (i.e. convolution of Fourier series corresponds to multiplication of represented periodic function and vice versa), but periodic functions cannot be convolved with the usual definition, since the involved integrals diverge. A possible way out is to define a periodic function on a bounded but periodic domain. To this end you can use the notion of a quotient space:
.
That is, each element in is an equivalence class of real numbers that share the same fractional part. Thus a function like is a representation of a 1-periodic function.
Calculating period
Consider a real waveform consisting of superimposed frequencies, expressed in a set as ratios to a fundamental frequency, f: F = [f f f ... f] where all non-zero elements ≥1 and at least one of the elements of the set is 1. To find the period, T, first find the least common denominator of all the elements in the set. Period can be found as T = . Consider that for a simple sinusoid, T = . Therefore, the LCD can be seen as a periodicity multiplier.
For set representing all notes of Western major scale: [1 ] the LCD is 24 therefore T = .
For set representing all notes of a major triad: [1 ] the LCD is 4 therefore T = .
For set representing all notes of a minor triad: [1 ] the LCD is 10 therefore T = .
If no least common denominator exists, for instance if one of the above elements were irrational, then the wave would not be periodic.
| Mathematics | Functions: General | null |
76885 | https://en.wikipedia.org/wiki/Guide%20dog | Guide dog | Guide dogs (colloquially known in the US as seeing-eye dogs) are assistance dogs trained to lead blind or visually impaired people around obstacles. Although dogs can be trained to navigate various obstacles, they are red–green colour blind and incapable of interpreting street signs. The human does the directing, based on skills acquired through previous mobility training. The handler might be likened to an aircraft's navigator, who must know how to get from one place to another, and the dog is the pilot, who gets them there safely. In several countries guide dogs, along with most other service and hearing dogs, are exempt from regulations against the presence of animals in places such as restaurants and public transportation.
History
| Technology | Agriculture, labor and economy | null |
76894 | https://en.wikipedia.org/wiki/Emu | Emu | The emu (; Dromaius novaehollandiae) is a species of flightless bird endemic to Australia, where it is the tallest native bird. It is the only extant member of the genus Dromaius and the third-tallest living bird after its African ratite relatives, the common ostrich and Somali ostrich. The emu's native ranges cover most of the Australian mainland. The Tasmanian, Kangaroo Island and King Island subspecies became extinct after the European settlement of Australia in 1788.
The emu has soft, brown feathers, a long neck, and long legs. It can grow up to in height. It is a robust bipedal runner that can travel great distances, and when necessary can sprint at . It is omnivorous and forages on a variety of plants and insects, and can go for weeks without eating. It drinks infrequently, but takes in copious amounts of fresh water when the opportunity arises.
Breeding takes place in May and June, and fighting among females for a mate is common. Females can mate several times and lay several clutches of eggs in one season. The male does the incubation; during this process he hardly eats or drinks and loses a significant amount of weight. The eggs hatch after around eight weeks, and the young are nurtured by their fathers. They reach full size after around six months, but can remain as a family unit until the next breeding season.
The emu is sufficiently common to be rated as a least-concern species by the International Union for Conservation of Nature. Despite this, some local populations are listed as endangered, with all the insular subspecies going extinct by the 1800s. Threats to their survival include egg predation by other animals (especially invasive species), roadkills and habitat fragmentation.
The emu is an important cultural icon of Australia, appearing on the coat of arms and various coinages. The bird features prominently in Indigenous Australian mythologies.
Etymology
The etymology of the common name "emu" is uncertain, but is thought to have come from an Arabic word for large bird that was later used by Portuguese explorers to describe the related cassowary in Australia and New Guinea. Another theory is that it comes from the word "ema", which is used in Portuguese to denote a large bird akin to an ostrich or crane. In Victoria, some terms for the emu were Barrimal in the Dja Dja Wurrung language, myoure in Gunai, and courn in Jardwadjali. The birds were known as murawung or birabayin to the local Eora and Darug inhabitants of the Sydney basin.
Taxonomy
History
Emus were first reported as having been seen by Europeans when explorers visited the western coast of Australia in 1696. This was during an expedition led by Dutch captain Willem de Vlamingh who was searching for survivors of a ship that had gone missing two years earlier. The birds were known on the eastern coast before 1788, when the first Europeans settled there. The birds were first mentioned under the name of the "New Holland cassowary" in Arthur Phillip's Voyage to Botany Bay, published in 1789 with the following description:
The species was named by ornithologist John Latham in 1790 based on a specimen from the Sydney area of Australia, a country which was known as New Holland at the time. He collaborated on Phillip's book and provided the first descriptions of, and names for, many Australian bird species; Dromaius comes from a Greek word meaning "racer" and novaehollandiae is the Latin term for New Holland, so the name can be rendered as "fast-footed New Hollander". In his original 1816 description of the emu, the French ornithologist Louis Pierre Vieillot used two generic names, first Dromiceius and later Dromaius. It has been a point of contention ever since as to which name should be used; the latter is more correctly formed, but the convention in taxonomy is that the first name given to an organism stands, unless it is clearly a typographical error. Most modern publications, including those of the Australian government, use Dromaius, with Dromiceius mentioned as an alternative spelling.
Systematics
The emu was long classified, with its closest relatives the cassowaries, in the family Casuariidae, part of the ratite order Struthioniformes. An alternate classification was proposed in 2014 by Mitchell et al., based on analysis of mitochondrial DNA. This splits off the Casuariidae into their own order, the Casuariformes, and includes only the cassowaries in the family Casuariidae, placing the emus in their own family, Dromaiidae. The cladogram shown below is from their study.
Two different Dromaius species were present in Australia at the time of European settlement, and one additional species is known from fossil remains. The insular dwarf emus, D. n. baudinianus and D. n. minor, originally present on Kangaroo Island and King Island respectively, both became extinct shortly after the arrival of Europeans. D. n. diemenensis, another insular dwarf emu from Tasmania, became extinct around 1865. The mainland subspecies, D. n. novaehollandiae, remains common. The population of these birds varies from decade to decade, largely being dependent on rainfall; in 2009, it was estimated that there were between 630,000 and 725,000 birds. Emus were introduced to Maria Island off Tasmania, and Kangaroo Island off the coast of South Australia, during the 20th century. The Maria Island population died out in the mid-1990s. The Kangaroo Island birds have successfully established a breeding population.
In 1912, the Australian ornithologist Gregory M. Mathews recognised three living subspecies of emu, D. n. novaehollandiae (Latham, 1790), D. n. woodwardi Mathews, 1912 and D. n. rothschildi Mathews, 1912. The Handbook of the Birds of the World, however, argues that the last two of these subspecies are invalid; natural variations in plumage colour and the nomadic nature of the species make it likely that there is a single race in mainland Australia. Examination of the DNA of the King Island emu shows this bird to be closely related to the mainland emu and hence best treated as a subspecies.
Description
The emu is the second tallest bird in the world, only being exceeded in height by the ostrich; the largest individuals can reach up to in height. Measured from the bill to the tail, emus range in length from , with males averaging and females averaging . Emus are the fourth or fifth heaviest living bird after the two species of ostrich and two larger species of cassowary, weighing slightly more on average than an emperor penguin. Adult emus weigh between , with an average of in males and females, respectively. Females are usually slightly larger than males and are substantially wider across the rump.
Although flightless, emus have vestigial wings, the wing chord measuring around , and each wing having a small claw at the tip. Emus flap their wings when running, perhaps as a means of stabilising themselves when moving fast. They have long necks and legs, and can run at speeds of due to their highly specialised pelvic limb musculature. Their feet have only three toes and a similarly reduced number of bones and associated foot muscles; emus are unique among birds in that their gastrocnemius muscles in the back of the lower legs have four bellies instead of the usual three. The pelvic limb muscles of emus contribute a similar proportion of the total body mass as do the flight muscles of flying birds. When walking, the emu takes strides of about , but at full gallop, a stride can be as long as . Its legs are devoid of feathers and underneath its feet are thick, cushioned pads. Like the cassowary, the emu has sharp claws on its toes which are its major defensive attribute, and are used in combat to inflict wounds on opponents by kicking. The toe and claw total in length. The bill is quite small, measuring , and is soft, being adapted for grazing. Emus have good eyesight and hearing, which allows them to detect threats at some distance.
The neck of the emu is pale blue and shows through its sparse feathers. They have grey-brown plumage of shaggy appearance; the shafts and the tips of the feathers are black. Solar radiation is absorbed by the tips, and the inner plumage insulates the skin. This prevents the birds from overheating, allowing them to be active during the heat of the day. A unique feature of the emu feather is the double rachis emerging from a single shaft. Both of the rachis have the same length, and the texture is variable; the area near the skin is rather furry, but the more distant ends resemble grass. The sexes are similar in appearance, although the male's penis can become visible when he urinates and defecates. The plumage varies in colour due to environmental factors, giving the bird a natural camouflage. Feathers of emus in more arid areas with red soils have a rufous tint while birds residing in damp conditions are generally darker in hue. The juvenile plumage develops at about three months and is blackish finely barred with brown, with the head and neck being especially dark. The facial feathers gradually thin to expose the bluish skin. The adult plumage has developed by about fifteen months.
The eyes of an emu are protected by nictitating membranes. These are translucent, secondary eyelids that move horizontally from the inside edge of the eye to the outside edge. They function as visors to protect the eyes from the dust that is prevalent in windy arid regions. Emus have a tracheal pouch, which becomes more prominent during the mating season. At more than in length, it is quite spacious; it has a thin wall, and an opening long.
Distribution and habitat
Once common on the east coast of Australia, emus are now uncommon there; by contrast, the development of agriculture and the provision of water for stock in the interior of the continent have increased the range of the emu in arid regions. Emus live in various habitats across Australia both inland and near the coast. They are most common in areas of savannah woodland and sclerophyll forest, and least common in heavily populated districts and arid areas with annual precipitation of less than . Emus predominantly travel in pairs, and while they can form large flocks, this is an atypical social behaviour that arises from the common need to move towards a new food source. Emus have been shown to travel long distances to reach abundant feeding areas. In Western Australia, emu movements follow a distinct seasonal pattern – north in summer and south in winter. On the east coast their wanderings seem to be more random and do not appear to follow a set pattern.
Behaviour and ecology
Emus are diurnal birds and spend their day foraging, preening their plumage with their beak, dust bathing and resting. They are generally gregarious birds apart from the breeding season, and while some forage, others remain vigilant to their mutual benefit. They are able to swim when necessary, although they rarely do so unless the area is flooded or they need to cross a river.
Emus begin to settle down at sunset and sleep during the night. They do not sleep continuously but rouse themselves several times during the night. When falling asleep, emus first squat on their tarsi and enter a drowsy state during which they are alert enough to react to stimuli and quickly return to a fully awakened state if disturbed. As they fall into deeper sleep, their neck droops closer to the body and the eyelids begin to close. If there are no disturbances, they fall into a deeper sleep after about twenty minutes. During this phase, the body is gradually lowered until it is touching the ground with the legs folded underneath. The beak is turned down so that the whole neck becomes S-shaped and folded onto itself. The feathers direct any rain downwards onto the ground. It has been suggested that the sleeping position is a type of camouflage, mimicking a small mound. Emus typically awake from deep sleep once every ninety minutes or so and stand upright to feed briefly or defecate. This period of wakefulness lasts for ten to twenty minutes, after which they return to slumber. Overall, an emu sleeps for around seven hours in each twenty-four-hour period. Young emus usually sleep with their neck flat and stretched forward along the ground surface.
The vocalisations of emus mostly consist of various booming and grunting sounds. The booming is created by the inflatable throat pouch; the pitch can be regulated by the bird and depends on the size of the aperture. Most of the booming is done by females; it is part of the courtship ritual, is used to announce the holding of territory and is issued as a threat to rivals. A high-intensity boom is audible away, while a low, more resonant call, produced during the breeding season, may at first attract mates and peaks while the male is incubating the eggs. Most of the grunting is done by males. It is used principally during the breeding season in territorial defence, as a threat to other males, during courtship and while the female is laying. Both sexes sometimes boom or grunt during threat displays or on encountering strange objects.
On very hot days, emus pant to maintain their body temperature. Their lungs work as evaporative coolers and, unlike some other species, the resulting low levels of carbon dioxide in the blood do not appear to cause alkalosis. For normal breathing in cooler weather, they have large, multifolded nasal passages. Cool air warms as it passes through into the lungs, extracting heat from the nasal region. On exhalation, the emu's cold nasal turbinates condense moisture back out of the air and absorb it for reuse. As with other ratites, the emu has great homeothermic ability, and can maintain this status from . The thermoneutral zone of emus lies between .
As with other ratites, emus have a relatively low basal metabolic rate compared to other types of birds. At , the metabolic rate of an emu sitting down is about 60% of that when standing, partly because the lack of feathers under the stomach leads to a higher rate of heat loss when standing from the exposed underbelly.
Diet
Emus forage in a diurnal pattern and eat a variety of native and introduced plant species. The diet depends on seasonal availability with such plants as Acacia, Casuarina and grasses being favoured. They also eat insects and other arthropods, including grasshoppers and crickets, beetles, cockroaches, ladybirds, bogong and cotton-boll moth larvae, ants, spiders and millipedes. This provides a large part of their protein requirements. In Western Australia, food preferences have been observed in travelling emus; they eat seeds from Acacia aneura until the rains arrive, after which they move on to fresh grass shoots and caterpillars; in winter they feed on the leaves and pods of Cassia and in spring, they consume grasshoppers and the fruit of Santalum acuminatum, a sort of quandong. They are also known to feed on wheat, and any fruit or other crops that they can access, easily climbing over high fences if necessary. Emus serve as an important agent for the dispersal of large viable seeds, which contributes to floral biodiversity. One undesirable effect of this occurred in Queensland in the early twentieth century when emus fed on the fruit of prickly pears in the outback. They defecated the seeds in various places as they moved around, and this led to a series of campaigns to hunt emus and prevent the seeds of the invasive cactus being spread. The cacti were eventually controlled by an introduced moth (Cactoblastis cactorum) whose larvae fed on the plant, one of the earliest examples of biological control. The δC of the emu's diet is reflected in the δC of the calcite of its egg shell.
Small stones are swallowed to assist in the grinding up and digestion of the plant material. Individual stones may weigh and the birds may have as much as in their gizzards at one time. They also eat charcoal, although the reason for this is unclear. Captive emus have been known to eat shards of glass, marbles, car keys, jewellery and nuts and bolts.
Emus drink infrequently but ingest large amounts when the opportunity arises. They typically drink once a day, first inspecting the water body and surrounding area in groups before kneeling down at the edge to drink. They prefer being on firm ground while drinking, rather than on rocks or mud, but if they sense danger, they often stand rather than kneel. If not disturbed, they may drink continuously for ten minutes. Due to the scarcity of water sources, emus are sometimes forced to go without water for several days. In the wild, they often share water holes with other animals such as kangaroos; they are wary and tend to wait for the other animals to leave before drinking.
Breeding
Emus form breeding pairs during the summer months of December and January and may remain together for about five months. During this time, they stay in an area a few kilometres in diameter and it is believed they find and defend territory within this area. Both males and females put on weight during the breeding season, with the female becoming slightly heavier at between . Mating usually takes place between April and June; the exact timing is determined by the climate as the birds nest during the coolest part of the year. During the breeding season, males experience hormonal changes, including an increase in luteinising hormone and testosterone levels, and their testicles double in size.
Males construct a rough nest in a semi-sheltered hollow on the ground, using bark, grass, sticks and leaves to line it. The nest is almost always a flat surface rather than a segment of a sphere, although in cold conditions the nest is taller, up to tall, and more spherical to provide some extra heat retention. When other material is lacking, the bird sometimes uses a spinifex tussock a metre or so across, despite the prickly nature of the foliage. The nest can be placed on open ground or near a shrub or rock. The nest is usually placed in an area where the emu has a clear view of its surroundings and can detect approaching predators. The nest can contain eggs from multiple emus the number is usually between 15 and 25 eggs.
Female emus court the males; the female's plumage darkens slightly and the small patches of bare, featherless skin just below the eyes and near the beak turn turquoise-blue. The colour of the male's plumage remains unchanged, although the bare patches of skin also turn light blue. When courting, females stride around, pulling their neck back while puffing out their feathers and emitting low, monosyllabic calls that have been compared to drum beats. This calling can occur when males are out of sight or more than away. Once the male's attention has been gained, the female circles her prospective mate at a distance of . As she does this, she looks at him by turning her neck, while at the same time keeping her rump facing towards him. If the male shows interest in the parading female, he will move closer; the female continues the courtship by shuffling further away but continuing to circle him.
If a male is interested, he will stretch his neck and erect his feathers, then bend over and peck at the ground. He will circle around and sidle up to the female, swaying his body and neck from side to side, and rubbing his breast against his partner's rump. Often the female will reject his advances with aggression, but if amenable, she signals acceptance by squatting down and raising her rump.
Females are more aggressive than males during the courtship period, often fighting for access to mates, with fights among females accounting for more than half the aggressive interactions during this period. If females court a male that already has a partner, the incumbent female will try to repel the competitor, usually by chasing and kicking. These interactions can be prolonged, lasting up to five hours, especially when the male being fought over is single and neither female has the advantage of incumbency. In these cases, the females typically intensify their calls and displays.
The sperm from a mating is stored by the female and can suffice to fertilise about six eggs. The pair mate every day or two, and every second or third day the female lays one of a clutch of five to fifteen very large, thick-shelled, green eggs. The shell is around thick, but rather thinner in northern regions according to indigenous Australians. The shell is substantially composed of calcite, and its δC is a function of the emu's diet. The eggs are on average and weigh between . The maternal investment in the egg is considerable, and the proportion of yolk to albumen, at about 50%, is greater than would be predicted for a precocial egg of this size. This probably relates to the long incubation period which means the developing chick must consume greater resources before hatching. The first verified occurrence of genetically identical avian twins was demonstrated in the emu. The egg surface is granulated and pale green. During the incubation period, the egg turns dark green, although if the egg never hatches, it will turn white from the bleaching effect of the sun.
The male becomes broody after his mate starts laying, and may begin to incubate the eggs before the clutch is complete. From this time on, he does not eat, drink, or defecate, and stands only to turn the eggs, which he does about ten times a day. He develops a brood patch, a bare area of wrinkled skin which is in intimate contact with the eggs. Over the course of the eight-week incubation period, he will lose a third of his weight and will survive on stored body fat and on any morning dew that he can reach from the nest. As with many other Australian birds, such as the superb fairywren, infidelity is the norm for emus, despite the initial pair bond: once the male starts brooding, the female usually wanders off, and may mate with other males and lay in multiple nests; thus, as many as half the chicks in a brood may not be fathered by the incubating male, or even by either parent, as emus also exhibit brood parasitism.
Some females stay and defend the nest until the chicks start hatching, but most leave the nesting area completely to nest again; in a good season, a female emu may nest three times. If the parents stay together during the incubation period, they will take turns standing guard over the eggs while the other drinks and feeds within earshot. If it perceives a threat during this period, it will lie down on top of the nest and try to blend in with the similar-looking surrounds, and suddenly stand up to confront and scare the other party if it comes close.
Incubation takes 56 days, and the male stops incubating the eggs shortly before they hatch. The temperature of the nest rises slightly during the eight-week period. Although the eggs are laid sequentially, they tend to hatch within two days of one another, as the eggs that were laid later experienced higher temperatures and developed more rapidly. During the process, the precocial emu chicks need to develop a capacity for thermoregulation. During incubation, the embryos are kept at a constant temperature but the chicks will need to be able to cope with varying external temperatures by the time they hatch.
Newly hatched chicks are active and can leave the nest within a few days of hatching. They stand about tall at first, weigh , and have distinctive brown and cream stripes for camouflage, which fade after three months or so. The male guards the growing chicks for up to seven months, teaching them how to find food. Chicks grow very quickly and are fully grown in five to six months; they may remain with their family group for another six months or so before they split up to breed in their second season. During their early life, the young emus are defended by their father, who adopts a belligerent stance towards other emus, including the mother. He does this by ruffling his feathers, emitting sharp grunts, and kicking his legs to drive off other animals. He can also bend his knees to crouch over smaller chicks to protect them. At night, he envelops his young with his feathers. As the young emus cannot travel far, the parents must choose an area with plentiful food in which to breed. In the wild, emus can live for upwards of 10 years but in captivity, they can live up to 20 years.
Predation
There are few native natural predators of adult emus still extant. Early in its species history it may have faced numerous terrestrial predators now extinct, including the giant lizard Megalania, the thylacine, and possibly other carnivorous marsupials, which may explain their seemingly well-developed ability to defend themselves from terrestrial predators. The main predator of emus today is the dingo, which was originally introduced by Aboriginals thousands of years ago from a stock of semi-domesticated wolves. Dingoes try to kill the emu by attacking the head. The emu typically tries to repel the dingo by jumping into the air and kicking or stamping the dingo on its way down. The emu jumps as the dingo barely has the capacity to jump high enough to threaten its neck, so a correctly timed leap to coincide with the dingo's lunge can keep its head and neck out of danger. Despite the potential prey-predator relationship, the presence of predaceous dingoes does not appear to heavily influence emu numbers, with other natural conditions just as likely to cause mortality. Wedge-tailed eagles are the only avian predator capable of attacking fully-grown emus, though are perhaps most likely to take small or young specimens. The eagles attack emus by swooping downwards rapidly and at high speed and aiming for the head and neck. In this case, the emu's jumping technique as employed against the dingo is not useful. The birds try to target the emu in the open ground so that it cannot hide behind obstacles. Under such circumstances, the emu runs in a chaotic manner and changes directions frequently to try to evade its attacker. While full-grown adults are rarely preyed upon, dingos, raptors, monitor lizards, introduced red foxes, feral and domestic dogs, and feral pigs occasionally feed on emu eggs or kill small chicks. Adult males fiercely defend their chicks from predators, especially dingos and foxes.
Parasites
Emus can suffer from both external and internal parasites, but under farmed conditions are more parasite-free than ostriches or rheas. External parasites include the louse Dahlemhornia asymmetrica and various other lice, ticks, mites and flies. Chicks sometimes suffer from intestinal tract infections caused by coccidian protozoa, and the nematode Trichostrongylus tenuis infects the emu as well as a wide range of other birds, causing haemorrhagic diarrhoea. Other nematodes are found in the trachea and bronchi; Syngamus trachea causing haemorrhagic tracheitis and Cyathostoma variegatum causing serious respiratory problems in juveniles.
Relationship with humans
Emus were used as a source of food by indigenous Australians and early European settlers. Emus are inquisitive birds and have been known to approach humans if they see unexpected movement of a limb or piece of clothing. In the wild, they may follow and observe people. Aboriginal Australians used a variety of techniques to catch the birds, including spearing them while they drank at waterholes, catching them in nets, and attracting them by imitating their calls or by arousing their curiosity with a ball of feathers and rags dangled from a tree. The pitchuri thornapple (Duboisia hopwoodii), or some similar poisonous plant, could be used to contaminate a waterhole, after which the disoriented emus were easy to catch. Another stratagem was for the hunter to use a skin as a disguise, and the birds could be lured into a camouflaged pit trap using rags or imitation calls. Aboriginal Australians only killed emus out of necessity, and frowned on anyone who hunted them for any other reason. Every part of the carcass had some use; the fat was harvested for its valuable, multiple-use oil, the bones were shaped into knives and tools, the feathers were used for body adornment and the tendons substituted for string.
The early European settlers killed emus to provide food and used their fat for fuelling lamps. They also tried to prevent them from interfering with farming or invading settlements in search of water during drought. An extreme example of this was the Emu War in Western Australia in 1932. Emus flocked to the Chandler and Walgoolan area during a dry spell, damaging rabbit fencing and devastating crops. An attempt to drive them off was mounted, with the army called in to dispatch them with machine guns; the emus largely avoided the hunters. Emus are large, powerful birds, and their legs are among the strongest of any animal and powerful enough to tear down metal fencing. The birds are very defensive of their young, and there have been two documented cases of humans being attacked by emus.
Economic value
In the areas in which it was endemic, the emu was an important source of meat to Aboriginal Australians. They used the fat as bush medicine and rubbed it into their skin. It served as a valuable lubricant, was used to oil wooden tools and utensils such as the coolamon, and was mixed with ochre to make the traditional paint for ceremonial body adornment. Their eggs were also foraged for food.
An example of how the emu was cooked comes from the Arrernte of Central Australia who called it Kere ankerre:
The birds were a food and fuel source for early European settlers, and are now farmed, in Australia and elsewhere, for their meat, oil and leather. Commercial emu farming started in Western Australia around 1970. The commercial industry in the country is based on stock bred in captivity, and all states except Tasmania have licensing requirements to protect wild emus. Outside Australia, emus are farmed on a large scale in North America, with about 1 million birds in the US, Peru, and China, and to a lesser extent in some other countries. Emus breed well in captivity, and are kept in large open pens to avoid the leg and digestive problems that arise from inactivity. They are typically fed on grain supplemented by grazing, and are slaughtered at 15 to 18 months.
The Salem district administration in India advised farmers in 2012 not to invest in the emu business which was being heavily promoted at the time; further investigation was needed to assess the profitability of farming the birds in India. In the United States, it was reported in 2013 that many ranchers had left the emu business; it was estimated that the number of growers had dropped from over five thousand in 1998 to one or two thousand in 2013. The remaining growers increasingly rely on sales of oil for their profit, although, leather, eggs, and meat are also sold.
Emus are farmed primarily for their meat, leather, feathers and oil, and 95% of the carcass can be used. Emu meat is a low-fat product (less than 1.5% fat), and is comparable to other lean meats. Most of the usable portions (the best cuts come from the thigh and the larger muscles of the drum or lower leg) are, like other poultry, dark meat; emu meat is considered for cooking purposes by the US Food and Drug Administration to be a red meat because its red colour and pH value approximate that of beef, but for inspection purposes it is considered to be poultry.
Emu fat is rendered to produce oil for cosmetics, dietary supplements, and therapeutic products. The oil is obtained from the subcutaneous and retroperitoneal fat; the macerated adipose tissue is heated and the liquefied fat is filtered to get a clear oil. This consists mainly of fatty acids of which oleic acid (42%), linoleic and palmitic acids (21% each) are the most prominent components. It also contains various anti-oxidants, notably carotenoids and flavones.
There is some evidence that the oil has anti-inflammatory properties; however, there have not yet been extensive tests, and the USDA regards pure emu oil as an unapproved drug and highlighted it in a 2009 article entitled "How to Spot Health Fraud". Nevertheless, the oil has been linked to the easing of gastrointestinal inflammation, and tests on rats have shown that it has a significant effect in treating arthritis and joint pain, more so than olive or fish oils. It has been scientifically shown to improve the rate of wound healing, but the mechanism responsible for this effect is not understood. A 2008 study has claimed that emu oil has a better anti-oxidative and anti-inflammatory potential than ostrich oil, and linked this to emu oil's higher proportion of unsaturated to saturated fatty acids. While there are no scientific studies showing that emu oil is effective in humans, it is marketed and promoted as a dietary supplement with a wide variety of claimed health benefits. Commercially marketed emu oil supplements are poorly standardised.
Emu leather has a distinctive patterned surface, due to a raised area around the feather follicles in the skin; the leather is used in such items as wallets, handbags, shoes and clothes, often in combination with other leathers. The feathers and eggs are used in decorative arts and crafts. In particular, emptied emu eggs have been engraved with portraits, similar to cameos, and scenes of Australian native animals. Mounted Emu eggs and emu-egg containers in the form of hundreds of goblets, inkstands and vases were produced in the second half of the nineteenth century, all richly embellished with images of Australian flora, fauna and indigenous people by travelling silversmiths, founders of a 'new Australian grammar of ornament'. They continued longstanding traditions that can be traced back to the European mounted ostrich eggs of the thirteenth century and Christian symbolism and notions of virginity, fertility, faith and strength. For a society of proud settlers who sought to bring culture and civilisation to their new world, the traditional ostrich-egg goblet, freed from its roots in a society dominated by court culture, was creatively made novel in the Australian colonies as forms and functions were invented to make the objects attractive to a new, broader audience. Significant designers Adolphus Blau, Julius Hogarth, Ernest Leviny, Julius Schomburgk, Johann Heinrich Steiner, Christian Quist, Joachim Matthias Wendt, William Edwards and others had the technical training on which to build flourishing businesses in a country rich in raw materials and a clientele hungry for old-world paraphernalia.
In addition to their use in farming, emus are sometimes kept as pets, though they require adequate space and food in order to live healthily. Emus were formerly subject to regulation in the United Kingdom under the Dangerous Wild Animals Act; however, a review of the act in 2007 led to changes that allow emus (alongside a number of other animals that were also regulated under the act) to be kept without a license, as they were no longer considered to be dangerous.
Cultural references
The emu has a prominent place in Australian Aboriginal mythology, including a creation myth of the Yuwaalaraay and other groups in New South Wales who say that the sun was made by throwing an emu's egg into the sky; the bird features in numerous aetiological stories told across a number of Aboriginal groups. One story from Western Australia holds that a man once annoyed a small bird, who responded by throwing a boomerang, severing the arms of the man and transforming him into a flightless emu. The Kurdaitcha man of Central Australia is said to wear sandals made of emu feathers to mask his footprints. Many Aboriginal language groups throughout Australia have a tradition that the dark dust lanes in the Milky Way represent a giant emu in the sky. Several of the Sydney rock engravings depict emus, and the birds are mimicked in Indigenous dances. Hunting emus, known as kari in the Kaurna language, features in the major Dreaming story of the Kaurna people of the Adelaide region about the ancestor hero Tjilbruke.
The emu is popularly but Unofficially considered as a faunal emblem – the national bird of Australia. It appears as a shield bearer on the Coat of arms of Australia with the red kangaroo, and as a part of the Arms also appears on the Australian 50-cent coin. It has featured on numerous Australian postage stamps, including a pre-federation New South Wales 100th Anniversary issue from 1888, which featured a 2 pence blue emu stamp, a 36-cent stamp released in 1986, and a $1.35 stamp released in 1994. The hats of the Australian Light Horse are decorated with emu feather plumes.
Trademarks of early Australian companies using the emu included Webbenderfer Bros frame mouldings (1891), Mac Robertson Chocolate and Cocoa (1893), Dyason and Son Emu Brand Cordial Sauce (1894), James Allard Pottery Wares (1906), and rope manufacturers G. Kinnear and Sons Pty. Ltd. still use it on some of their products.
There are around six hundred gazetted places in Australia with "emu" in their title, including mountains, lakes, hills, plains, creeks and waterholes. During the 19th and 20th centuries, many Australian companies and household products were named after the bird. In Western Australia, Emu beer has been produced since the early 20th century and the Swan Brewery continues to produce a range of beers branded as "Emu". The quarterly peer-reviewed journal of the Royal Australasian Ornithologists Union, also known as Birds Australia, is entitled Emu: Austral Ornithology.
The comedian Rod Hull featured a wayward emu puppet in his act for many years and the bird returned to the small screen in the hands of his son Toby after the puppeteer's death in 1999. In 2019, American insurance company Liberty Mutual launched an advertising campaign that features LiMu Emu, a CGI-rendered emu.
Another popular Emu on social media is Emmanuel, a resident of Knuckle Bump Farms in south Florida. Taylor Blake, the farm's owner, since 2013 has recorded video shorts explaining aspects of the farm and is often interrupted as Emmanuel the Emu photobombs her videos earning constant rebukes; the term "Emmanuel don't do it!" has become popular on social media.
Status and conservation
In John Gould's Handbook to the Birds of Australia, first published in 1865, he lamented the loss of the emu from Tasmania, where it had become rare and has since become extinct; he noted that emus were no longer common in the vicinity of Sydney and proposed that the species be given protected status. In the 1930s, emu killings in Western Australia peaked at 57,000, and culls were also mounted in Queensland during this period due to rampant crop damage. In the 1960s, bounties were still being paid in Western Australia for killing emus, but since then, wild emus have been granted formal protection under the Environment Protection and Biodiversity Conservation Act 1999. Their occurrence range is between , and a 1992 census suggested that their total population was between 630,000 and 725,000. Their population trend is thought to be stable and the International Union for Conservation of Nature assesses their conservation status as being of least concern. The isolated emu population of the New South Wales North Coast Bioregion and Port Stephens is listed as endangered by the New South Wales Government.
Although the population of emus on mainland Australia is thought to be higher now than it was before European settlement, some local populations are at risk of extinction. The threats faced by emus include the clearance and fragmentation of areas of suitable habitat, deliberate slaughter, collisions with vehicles and predation of the eggs and young.
| Biology and health sciences | Ratites | null |
76902 | https://en.wikipedia.org/wiki/Nicotiana | Nicotiana | Nicotiana () is a genus of herbaceous plants and shrubs in the family Solanaceae, that is indigenous to the Americas, Australia, Southwestern Africa and the South Pacific. Various Nicotiana species, commonly referred to as tobacco plants, are cultivated as ornamental garden plants. N. tabacum is grown worldwide for the cultivation of tobacco leaves used for manufacturing and producing tobacco products, including cigars, cigarillos, cigarettes, chewing tobacco, dipping tobacco, snuff, and snus.
Taxonomy
Species
The 79 accepted and known species include:
Nicotiana acuminata (Graham) Hook. – manyflower tobacco or many-flowered tobacco
Nicotiana africana Merxm.
Nicotiana alata Link & Otto – jasmine tobacco, sweet tobacco, winged tobacco, Persian tobacco, tanbaku (in Persian)
Nicotiana attenuata Torrey ex S. Watson – coyote tobacco
Nicotiana benthamiana Domin – benth, benthi
Nicotiana clevelandii A. Gray – Cleveland's tobacco
Nicotiana glauca Graham – tree tobacco, Brazilian tree tobacco, shrub tobacco, wild tobacco, tobacco plant, tobacco bush, tobacco tree, mustard tree
Nicotiana glutinosa L.
Nicotiana langsdorffii Weinm. – Langsdorff's tobacco
Nicotiana longiflora Cav. – longflower tobacco or long-flowered tobacco
Nicotiana occidentalis H.-M. Wheeler – native tobacco
Nicotiana obtusifolia M. Martens & Galeotti – desert tobacco, punche, "tabaquillo"
Nicotiana otophora Griseb.
Nicotiana plumbaginifolia Viv. – Tex-Mex tobacco
Nicotiana quadrivalvis Pursh – Indian tobacco
Nicotiana rustica L. – Aztec tobacco, strong tobacco, mapacho
Nicotiana suaveolens Lehm. – Australian tobacco
Nicotiana sylvestris Speg. & Comes – woodland tobacco, flowering tobacco, South American tobacco
Nicotiana tabacum L. – common tobacco, domesticated tobacco, cultivated tobacco, commercial tobacco (grown for the production of cigars, cigarillos, cigarettes, chewing tobacco, dipping tobacco, snuff, snus, etc.)
Nicotiana tomentosiformis Goodsp.
Manmade hybrids
Nicotiana × didepta – N. forsteri × N. tabacum
Nicotiana × digluta – N. glutinosa × N. tabacum
Nicotiana × sanderae Hort. ex Wats. – N. alata × N. forgetiana
Formerly placed here
Petunia axillaris (Lam.) Britton et al. (as N. axillaris Lam.) – large white petunia, wild white petunia, white moon petunia
Etymology
The genus Nicotiana (from which the word nicotine is derived) was named in honor of Jean Nicot, French ambassador to Portugal, who in 1559 sent samples as a medicine to the court of Catherine de' Medici.
Ecology
Despite containing enough nicotine and/or other compounds such as germacrene and anabasine and other piperidine alkaloids (varying between species) to deter most herbivores, a number of such animals have evolved the ability to feed on Nicotiana species without being harmed. Nonetheless, tobacco is unpalatable to many species and therefore some tobacco plants (mainly tree tobacco (N. glauca)) have become established as invasive species in some places.
In the 19th century, young tobacco plantings came under increasing attack from flea beetles (particularly the potato flea beetle (Epitrix cucumeris) and/or Epitrix pubescens), causing the destruction of half the United States tobacco crop in 1876. In the years afterward, many experiments were attempted and discussed to control the potato flea beetle. By 1880, it was discovered that covering young plants with a frame covered with thin fabric (instead of with branches, as had previously been used for frost control) would effectively protect the plants from the beetle. This practice spread until it became ubiquitous in the 1890s.
Tobacco, alongside its related products, can be infested by parasites such as the tobacco beetle (Lasioderma serricorne) and the tobacco moth (Ephestia elutella), which are the most widespread and damaging pests in the tobacco industry. Infestation can range from the tobacco cultivated in the fields to the leaves used for manufacturing cigars, cigarillos, cigarettes, chewing tobacco, dipping tobacco, snuff, snus, etc. Both the grubs of Lasioderma serricorne and the caterpillars of Ephestia elutella are considered major pests.
Other moths whose caterpillars feed on Nicotiana include:
Black cutworm, greasy cutworm, or floodplain cutworm (as a caterpillar), dark sword-grass or ipsilon dart (as a moth) (Agrotis ipsilon)
Turnip moth (Agrotis segetum)
Mouse moth (Amphipyra tragopoginis)
Clover cutworm (as a caterpillar), nutmeg (as a moth) (Hadula trifolii or Anarta trifolii)
Endoclita excrescens
Hawaiian tobacco hornworm or Hawaiian tomato hornworm (as a caterpillar), Blackburn's sphinx moth (as a moth) (Manduca blackburni)
Tobacco hornworm or Goliath worm (as a caterpillar), tobacco hawkmoth or Carolina sphinx moth (as a moth) (Manduca sexta)
Tomato hornworm (as a caterpillar), five-spotted hawkmoth (as a moth) (Manduca quinquemaculata)
Cabbage moth (Mamestra brassicae)
Angle shades (Phlogophora meticulosa)
Setaceous Hebrew character (Xestia c-nigrum)
Cabbage looper (Trichoplusia ni)
Fall armyworm (Spodoptera frugiperda)
Tobacco spitworm (as a caterpillar), potato tuber moth (as a moth) (Phthorimaea operculella)
South American tomato pinworm, tomato pinworm or tomato leafminer (as a caterpillar), South American tomato moth (as a moth) (Tuta absoluta)
Eggplant leafroller moth or nightshade leaftier (Lineodes integra)
Eggplant webworm moth (Rhectocraspeda periusalis)
These are mainly Noctuidae, but they also comprise Sphingidae, Gelechiidae, and Crambidae.
Cultivation
Several species of Nicotiana, such as N. sylvestris, N. alata 'Lime Green' and N. langsdorffii are grown as ornamental plants, often under the name of flowering tobacco. They are popular vespertines (evening bloomers); their sweet-smelling flowers opening in the evening to be visited by hawkmoths and other pollinators. In temperate climates, they behave as annuals (hardiness 9a-11).
The hybrid cultivar 'Lime Green' has gained the Royal Horticultural Society's Award of Garden Merit.
Garden varieties are derived from N. alata (e.g., the 'Niki' and 'Saratoga' series) and more recently from Nicotiana × sanderae (e.g., the 'Perfume' and 'Domino' series).
The tobacco budworm (Chloridea virescens) has proved to be a massive "pest" of many species in the genus, and has resisted many attempts at management.
| Biology and health sciences | Solanales | Plants |
76944 | https://en.wikipedia.org/wiki/Perpendicular | Perpendicular | In geometry, two geometric objects are perpendicular if they intersect at right angles, i.e. at an angle of 90 degrees or π/2 radians. The condition of perpendicularity may be represented graphically using the perpendicular symbol, ⟂. Perpendicular intersections can happen between two lines (or two line segments), between a line and a plane, and between two planes.
Perpendicular is also used as a noun: a perpendicular is a line which is perpendicular to a given line or plane.
Perpendicularity is one particular instance of the more general mathematical concept of orthogonality; perpendicularity is the orthogonality of classical geometric objects. Thus, in advanced mathematics, the word "perpendicular" is sometimes used to describe much more complicated geometric orthogonality conditions, such as that between a surface and its normal vector.
A line is said to be perpendicular to another line if the two lines intersect at a right angle. Explicitly, a first line is perpendicular to a second line if (1) the two lines meet; and (2) at the point of intersection the straight angle on one side of the first line is cut by the second line into two congruent angles. Perpendicularity can be shown to be symmetric, meaning if a first line is perpendicular to a second line, then the second line is also perpendicular to the first. For this reason, we may speak of two lines as being perpendicular (to each other) without specifying an order. A great example of perpendicularity can be seen in any compass, note the cardinal points; North, East, South, West (NESW)
The line N-S is perpendicular to the line W-E and the angles N-E, E-S, S-W and W-N are all 90° to one another.
Perpendicularity easily extends to segments and rays. For example, a line segment is perpendicular to a line segment if, when each is extended in both directions to form an infinite line, these two resulting lines are perpendicular in the sense above. In symbols, means line segment AB is perpendicular to line segment CD.
A line is said to be perpendicular to a plane if it is perpendicular to every line in the plane that it intersects. This definition depends on the definition of perpendicularity between lines.
Two planes in space are said to be perpendicular if the dihedral angle at which they meet is a right angle.
Foot of a perpendicular
The word foot is frequently used in connection with perpendiculars. This usage is exemplified in the top diagram, above, and its caption. The diagram can be in any orientation. The foot is not necessarily at the bottom.
More precisely, let be a point and a line. If is the point of intersection of and the unique line through that is perpendicular to , then is called the foot of this perpendicular through .
Construction of the perpendicular
To make the perpendicular to the line AB through the point P using compass-and-straightedge construction, proceed as follows (see figure left):
Step 1 (red): construct a circle with center at P to create points A' and B' on the line AB, which are equidistant from P.
Step 2 (green): construct circles centered at A' and B' having equal radius. Let Q and P be the points of intersection of these two circles.
Step 3 (blue): connect Q and P to construct the desired perpendicular PQ.
To prove that the PQ is perpendicular to AB, use the SSS congruence theorem for QPA' and QPB' to conclude that angles OPA' and OPB' are equal. Then use the SAS congruence theorem for triangles OPA' and OPB' to conclude that angles POA and POB are equal.
To make the perpendicular to the line g at or through the point P using Thales's theorem, see the animation at right.
The Pythagorean theorem can be used as the basis of methods of constructing right angles. For example, by counting links, three pieces of chain can be made with lengths in the ratio 3:4:5. These can be laid out to form a triangle, which will have a right angle opposite its longest side. This method is useful for laying out gardens and fields, where the dimensions are large, and great accuracy is not needed. The chains can be used repeatedly whenever required.
In relationship to parallel lines
If two lines (a and b) are both perpendicular to a third line (c), all of the angles formed along the third line are right angles. Therefore, in Euclidean geometry, any two lines that are both perpendicular to a third line are parallel to each other, because of the parallel postulate. Conversely, if one line is perpendicular to a second line, it is also perpendicular to any line parallel to that second line.
In the figure at the right, all of the orange-shaded angles are congruent to each other and all of the green-shaded angles are congruent to each other, because vertical angles are congruent and alternate interior angles formed by a transversal cutting parallel lines are congruent. Therefore, if lines a and b are parallel, any of the following conclusions leads to all of the others:
One of the angles in the diagram is a right angle.
One of the orange-shaded angles is congruent to one of the green-shaded angles.
Line c is perpendicular to line a.
Line c is perpendicular to line b.
All four angles are equal.
In computing distances
Graph of functions
In the two-dimensional plane, right angles can be formed by two intersected lines if the product of their slopes equals −1. Thus for two linear functions and , the graphs of the functions will be perpendicular if
The dot product of vectors can be also used to obtain the same result: First, shift coordinates so that the origin is situated where the lines cross. Then define two displacements along each line, , for Now, use the fact that the inner product vanishes for perpendicular vectors:
(unless or vanishes.)
Both proofs are valid for horizontal and vertical lines to the extent that we can let one slope be , and take the limit that If one slope goes to zero, the other goes to infinity.
In circles and other conics
Circles
Each diameter of a circle is perpendicular to the tangent line to that circle at the point where the diameter intersects the circle.
A line segment through a circle's center bisecting a chord is perpendicular to the chord.
If the intersection of any two perpendicular chords divides one chord into lengths a and b and divides the other chord into lengths c and d, then equals the square of the diameter.
The sum of the squared lengths of any two perpendicular chords intersecting at a given point is the same as that of any other two perpendicular chords intersecting at the same point, and is given by 8r2 – 4p2 (where r is the circle's radius and p is the distance from the center point to the point of intersection).
Thales' theorem states that two lines both through the same point on a circle but going through opposite endpoints of a diameter are perpendicular. This is equivalent to saying that any diameter of a circle subtends a right angle at any point on the circle, except the two endpoints of the diameter.
Ellipses
The major and minor axes of an ellipse are perpendicular to each other and to the tangent lines to the ellipse at the points where the axes intersect the ellipse.
The major axis of an ellipse is perpendicular to the directrix and to each latus rectum.
Parabolas
In a parabola, the axis of symmetry is perpendicular to each of the latus rectum, the directrix, and the tangent line at the point where the axis intersects the parabola.
From a point on the tangent line to a parabola's vertex, the other tangent line to the parabola is perpendicular to the line from that point through the parabola's focus.
The orthoptic property of a parabola is that If two tangents to the parabola are perpendicular to each other, then they intersect on the directrix. Conversely, two tangents which intersect on the directrix are perpendicular. This implies that, seen from any point on its directrix, any parabola subtends a right angle.
Hyperbolas
The transverse axis of a hyperbola is perpendicular to the conjugate axis and to each directrix.
The product of the perpendicular distances from a point P on a hyperbola or on its conjugate hyperbola to the asymptotes is a constant independent of the location of P.
A rectangular hyperbola has asymptotes that are perpendicular to each other. It has an eccentricity equal to
In polygons
Triangles
The legs of a right triangle are perpendicular to each other.
The altitudes of a triangle are perpendicular to their respective bases. The perpendicular bisectors of the sides also play a prominent role in triangle geometry.
The Euler line of an isosceles triangle is perpendicular to the triangle's base.
The Droz-Farny line theorem concerns a property of two perpendicular lines intersecting at a triangle's orthocenter.
Harcourt's theorem concerns the relationship of line segments through a vertex and perpendicular to any line tangent to the triangle's incircle.
Quadrilaterals
In a square or other rectangle, all pairs of adjacent sides are perpendicular. A right trapezoid is a trapezoid that has two pairs of adjacent sides that are perpendicular.
Each of the four maltitudes of a quadrilateral is a perpendicular to a side through the midpoint of the opposite side.
An orthodiagonal quadrilateral is a quadrilateral whose diagonals are perpendicular. These include the square, the rhombus, and the kite. By Brahmagupta's theorem, in an orthodiagonal quadrilateral that is also cyclic, a line through the midpoint of one side and through the intersection point of the diagonals is perpendicular to the opposite side.
By van Aubel's theorem, if squares are constructed externally on the sides of a quadrilateral, the line segments connecting the centers of opposite squares are perpendicular and equal in length.
Lines in three dimensions
Up to three lines in three-dimensional space can be pairwise perpendicular, as exemplified by the x, y, and z axes of a three-dimensional Cartesian coordinate system.
| Mathematics | Two-dimensional space | null |
76956 | https://en.wikipedia.org/wiki/Right%20angle | Right angle | In geometry and trigonometry, a right angle is an angle of exactly 90 degrees or radians corresponding to a quarter turn. If a ray is placed so that its endpoint is on a line and the adjacent angles are equal, then they are right angles. The term is a calque of Latin angulus rectus; here rectus means "upright", referring to the vertical perpendicular to a horizontal base line.
Closely related and important geometrical concepts are perpendicular lines, meaning lines that form right angles at their point of intersection, and orthogonality, which is the property of forming right angles, usually applied to vectors. The presence of a right angle in a triangle is the defining factor for right triangles, making the right angle basic to trigonometry.
Etymology
The meaning of right in right angle possibly refers to the Latin adjective rectus 'erect, straight, upright, perpendicular'. A Greek equivalent is orthos 'straight; perpendicular' (see orthogonality).
In elementary geometry
A rectangle is a quadrilateral with four right angles. A square has four right angles, in addition to equal-length sides.
The Pythagorean theorem states how to determine when a triangle is a right triangle.
Symbols
In Unicode, the symbol for a right angle is . It should not be confused with the similarly shaped symbol . Related symbols are , , and .
In diagrams, the fact that an angle is a right angle is usually expressed by adding a small right angle that forms a square with the angle in the diagram, as seen in the diagram of a right triangle (in British English, a right-angled triangle) to the right. The symbol for a measured angle, an arc, with a dot, is used in some European countries, including German-speaking countries and Poland, as an alternative symbol for a right angle.
Euclid
Right angles are fundamental in Euclid's Elements. They are defined in Book 1, definition 10, which also defines perpendicular lines. Definition 10 does not use numerical degree measurements but rather touches at the very heart of what a right angle is, namely two straight lines intersecting to form two equal and adjacent angles. The straight lines which form right angles are called perpendicular. Euclid uses right angles in definitions 11 and 12 to define acute angles (those smaller than a right angle) and obtuse angles (those greater than a right angle). Two angles are called complementary if their sum is a right angle.
Book 1 Postulate 4 states that all right angles are equal, which allows Euclid to use a right angle as a unit to measure other angles with. Euclid's commentator Proclus gave a proof of this postulate using the previous postulates, but it may be argued that this proof makes use of some hidden assumptions. Saccheri gave a proof as well but using a more explicit assumption. In Hilbert's axiomatization of geometry this statement is given as a theorem, but only after much groundwork. One may argue that, even if postulate 4 can be proven from the preceding ones, in the order that Euclid presents his material it is necessary to include it since without it postulate 5, which uses the right angle as a unit of measure, makes no sense.
Conversion to other units
A right angle may be expressed in different units:
turn
90° (degrees)
radians
100 grad (also called grade, gradian, or gon)
8 points (of a 32-point compass rose)
6 hours (astronomical hour angle)
Rule of 3-4-5
Throughout history, carpenters and masons have known a quick way to confirm if an angle is a true right angle. It is based on the Pythagorean triple and the rule of 3-4-5. From the angle in question, running a straight line along one side exactly three units in length, and along the second side exactly four units in length, will create a hypotenuse (the longer line opposite the right angle that connects the two measured endpoints) of exactly five units in length.
Thales' theorem
Thales' theorem states that an angle inscribed in a semicircle (with a vertex on the semicircle and its defining rays going through the endpoints of the semicircle) is a right angle.
Two application examples in which the right angle and the Thales' theorem are included (see animations).
Generalizations
The solid angle subtended by an octant of a sphere (the spherical triangle with three right angles) equals /2 sr.
| Mathematics | Geometry | null |
76988 | https://en.wikipedia.org/wiki/Electrocardiography | Electrocardiography | Electrocardiography is the process of producing an electrocardiogram (ECG or EKG), a recording of the heart's electrical activity through repeated cardiac cycles. It is an electrogram of the heart which is a graph of voltage versus time of the electrical activity of the heart using electrodes placed on the skin. These electrodes detect the small electrical changes that are a consequence of cardiac muscle depolarization followed by repolarization during each cardiac cycle (heartbeat). Changes in the normal ECG pattern occur in numerous cardiac abnormalities, including:
Cardiac rhythm disturbances, such as atrial fibrillation and ventricular tachycardia;
Inadequate coronary artery blood flow, such as myocardial ischemia and myocardial infarction;
and electrolyte disturbances, such as hypokalemia.
Traditionally, "ECG" usually means a 12-lead ECG taken while lying down as discussed below.
However, other devices can record the electrical activity of the heart such as a Holter monitor but also some models of smartwatch are capable of recording an ECG.
ECG signals can be recorded in other contexts with other devices.
In a conventional 12-lead ECG, ten electrodes are placed on the patient's limbs and on the surface of the chest. The overall magnitude of the heart's electrical potential is then measured from twelve different angles ("leads") and is recorded over a period of time (usually ten seconds). In this way, the overall magnitude and direction of the heart's electrical depolarization is captured at each moment throughout the cardiac cycle.
There are three main components to an ECG:
The P wave, which represents depolarization of the atria.
The QRS complex, which represents depolarization of the ventricles.
The T wave, which represents repolarization of the ventricles.
During each heartbeat, a healthy heart has an orderly progression of depolarization that starts with pacemaker cells in the sinoatrial node, spreads throughout the atrium, and passes through the atrioventricular node down into the bundle of His and into the Purkinje fibers, spreading down and to the left throughout the ventricles. This orderly pattern of depolarization gives rise to the characteristic ECG tracing. To the trained clinician, an ECG conveys a large amount of information about the structure of the heart and the function of its electrical conduction system. Among other things, an ECG can be used to measure the rate and rhythm of heartbeats, the size and position of the heart chambers, the presence of any damage to the heart's muscle cells or conduction system, the effects of heart drugs, and the function of implanted pacemakers.
Medical uses
The overall goal of performing an ECG is to obtain information about the electrical functioning of the heart. Medical uses for this information are varied and often need to be combined with knowledge of the structure of the heart and physical examination signs to be interpreted. Some indications for performing an ECG include the following:
Chest pain or suspected myocardial infarction (heart attack), such as ST elevated myocardial infarction (STEMI) or non-ST elevated myocardial infarction (NSTEMI)
Symptoms such as shortness of breath, murmurs, fainting, seizures, funny turns, or arrhythmias including new onset palpitations or monitoring of known cardiac arrhythmias
Medication monitoring (e.g., drug-induced QT prolongation, digoxin toxicity) and management of overdose (e.g., tricyclic overdose)
Electrolyte abnormalities, such as hyperkalemia
Perioperative monitoring in which any form of anesthesia is involved (e.g., monitored anesthesia care, general anesthesia). This includes preoperative assessment and intraoperative and postoperative monitoring.
Cardiac stress testing
Computed tomography angiography (CTA) and magnetic resonance angiography (MRA) of the heart (ECG is used to "gate" the scanning so that the anatomical position of the heart is steady)
Clinical cardiac electrophysiology, in which a catheter is inserted through the femoral vein and can have several electrodes along its length to record the direction of electrical activity from within the heart.
ECGs can be recorded as short intermittent tracings or continuous ECG monitoring. Continuous monitoring is used for critically ill patients, patients undergoing general anesthesia, and patients who have an infrequently occurring cardiac arrhythmia that would unlikely be seen on a conventional ten-second ECG. Continuous monitoring can be conducted by using Holter monitors, internal and external defibrillators and pacemakers, and/or biotelemetry.
Screening
For adults, evidence does not support the use of ECGs among those without symptoms or at low risk of cardiovascular disease as an effort for prevention. This is because an ECG may falsely indicate the existence of a problem, leading to misdiagnosis, the recommendation of invasive procedures, and overtreatment. However, persons employed in certain critical occupations, such as aircraft pilots, may be required to have an ECG as part of their routine health evaluations. Hypertrophic cardiomyopathy screening may also be considered in adolescents as part of a sports physical out of concern for sudden cardiac death.
Electrocardiograph machines
Electrocardiograms are recorded by machines that consist of a set of electrodes connected to a central unit. Early ECG machines were constructed with analog electronics, where the signal drove a motor to print out the signal onto paper. Today, electrocardiographs use analog-to-digital converters to convert the electrical activity of the heart to a digital signal. Many ECG machines are now portable and commonly include a screen, keyboard, and printer on a small wheeled cart. Recent advancements in electrocardiography include developing even smaller devices for inclusion in fitness trackers and smart watches. These smaller devices often rely on only two electrodes to deliver a single lead I. Portable twelve-lead devices powered by batteries are also available.
Recording an ECG is a safe and painless procedure. The machines are powered by mains power but they are designed with several safety features including an earthed (ground) lead.
Other features include:
Defibrillation protection: any ECG used in healthcare may be attached to a person who requires defibrillation and the ECG needs to protect itself from this source of energy.
Electrostatic discharge is similar to defibrillation discharge and requires voltage protection up to 18,000 volts.
Additionally, circuitry called the right leg driver can be used to reduce common-mode interference (typically the 50 or 60 Hz mains power).
ECG voltages measured across the body are very small. This low voltage necessitates a low noise circuit, instrumentation amplifiers, and electromagnetic shielding.
Simultaneous lead recordings: earlier designs recorded each lead sequentially, but current models record multiple leads simultaneously.
Most modern ECG machines include automated interpretation algorithms. This analysis calculates features such as the PR interval, QT interval, corrected QT (QTc) interval, PR axis, QRS axis, rhythm and more. The results from these automated algorithms are considered "preliminary" until verified and/or modified by expert interpretation. Despite recent advances, computer misinterpretation remains a significant problem and can result in clinical mismanagement.
Cardiac monitors
Besides the standard electrocardiograph machine, there are other devices capable of recording ECG signals. Portable devices have existed since the Holter monitor was introduced in 1962.
Traditionally, these monitors have used electrodes with patches on the skin to record the ECG, but new devices can stick to the chest as a single patch without need for wires, developed by Zio (Zio XT), TZ Medical (Trident), Philips (BioTel) and BardyDx (CAM) among many others.
Implantable devices such as the artificial cardiac pacemaker and implantable cardioverter-defibrillator are capable of measuring a "far field" signal between the leads in the heart and the implanted battery/generator that resembles an ECG signal (technically, the signal recorded in the heart is called an electrogram, which is interpreted differently).
Advancement of the Holter monitor became the implantable loop recorder that performs the same function but in an implantable device with batteries that last on the order of years.
Additionally, there are available various Arduino kits with ECG sensor modules and smartwatch devices that are capable of recording an ECG signal as well, such as with the 4th generation Apple Watch, Samsung Galaxy Watch 4 and newer devices.
Electrodes and leads
Electrodes are the actual conductive pads attached to the body surface. Any pair of electrodes can measure the electrical potential difference between the two corresponding locations of attachment. Such a pair forms a lead. However, "leads" can also be formed between a physical electrode and a virtual electrode, known as Wilson's central terminal (WCT), whose potential is defined as the average potential measured by three limb electrodes that are attached to the right arm, the left arm, and the left foot, respectively.
Commonly, 10 electrodes attached to the body are used to form 12 ECG leads, with each lead measuring a specific electrical potential difference (as listed in the table below).
Electrodes applied to patient's body
Leads are broken down into three types: limb; augmented limb; and precordial or chest. The 12-lead ECG has a total of three limb leads and three augmented limb leads arranged like spokes of a wheel in the coronal plane (vertical), and six precordial leads or chest leads that lie on the perpendicular transverse plane (horizontal).
Leads should be placed in standard positions. Exceptions due to emergency or other issues should be recorded to avoid erroneous analysis.
The 12 standard ECG leads are listed below. All leads are effectively bipolar, with one positive and one negative electrode; the term "unipolar" is not useful.
Two types of electrodes in common use are a flat paper-thin sticker and a self-adhesive circular pad.
The former are typically used in a single ECG recording while the latter are for continuous recordings as they stick longer.
Each electrode consists of an electrically conductive electrolyte gel and a silver/silver chloride conductor.
The gel typically contains potassium chloride – sometimes silver chloride as well – to permit electron conduction from the skin to the wire and to the electrocardiogram.
The common virtual electrode, known as Wilson's central terminal (VW), is produced by averaging the measurements from the electrodes RA, LA, and LL to give an average potential of the body:
In a 12-lead ECG, all leads except the limb leads are assumed to be unipolar (aVR, aVL, aVF, V1, V2, V3, V4, V5, and V6).
The measurement of a voltage requires two contacts and so, electrically, the unipolar leads are measured from the common lead (negative) and the unipolar lead (positive).
This averaging for the common lead and the abstract unipolar lead concept makes for a more challenging understanding and is complicated by sloppy usage of "lead" and "electrode".
In fact, instead of being a constant reference, VW has a value that fluctuates throughout the heart cycle. It also does not truly represent the center-of-heart potential due to the body parts the signals travel through. Because voltage is by definition a bipolar measurement between two points, describing an electrocardiographic lead as "unipolar" makes little sense electrically and should be avoided. The American Heart Association states "All leads are effectively 'bipolar,' and the term 'unipolar' in description of the augmented limb leads and the precordial leads lacks precision."
Limb leads
Leads I, II and III are called the limb leads. The electrodes that form these signals are located on the limbs – one on each arm and one on the left leg. The limb leads form the points of what is known as Einthoven's triangle.
Lead I is the voltage between the (positive) left arm (LA) electrode and right arm (RA) electrode:
Lead II is the voltage between the (positive) left leg (LL) electrode and the right arm (RA) electrode:
Lead III is the voltage between the (positive) left leg (LL) electrode and the left arm (LA) electrode:
Augmented limb leads
Leads aVR, aVL, and aVF are the augmented limb leads. They are derived from the same three electrodes as leads I, II, and III, but they use Goldberger's central terminal as their negative pole. Goldberger's central terminal is a combination of inputs from two limb electrodes, with a different combination for each augmented lead. It is referred to immediately below as "the negative pole".
Lead augmented vector right (aVR) has the positive electrode on the right arm. The negative pole is a combination of the left arm electrode and the left leg electrode:
Lead augmented vector left (aVL) has the positive electrode on the left arm. The negative pole is a combination of the right arm electrode and the left leg electrode:
Lead augmented vector foot (aVF) has the positive electrode on the left leg. The negative pole is a combination of the right arm electrode and the left arm electrode:
Together with leads I, II, and III, augmented limb leads aVR, aVL, and aVF form the basis of the hexaxial reference system, which is used to calculate the heart's electrical axis in the frontal plane.
Older versions of the nodes (VR, VL, VF) use Wilson's central terminal as the negative pole, but the amplitude is too small for the thick lines of old ECG machines. The Goldberger terminals scale up (augments) the Wilson results by 50%, at the cost of sacrificing physical correctness by not having the same negative pole for all three.
Precordial leads
The precordial leads lie in the transverse (horizontal) plane, perpendicular to the other six leads. The six precordial electrodes act as the positive poles for the six corresponding precordial leads: (V1, V2, V3, V4, V5, and V6). Wilson's central terminal is used as the negative pole. Recently, unipolar precordial leads have been used to create bipolar precordial leads that explore the right to left axis in the horizontal plane.
Specialized leads
Additional electrodes may rarely be placed to generate other leads for specific diagnostic purposes. Right-sided precordial leads may be used to better study pathology of the right ventricle or for dextrocardia (and are denoted with an R (e.g., V5R). Posterior leads (V7 to V9) may be used to demonstrate the presence of a posterior myocardial infarction. The Lewis lead or S5-lead (requiring an electrode at the right sternal border in the second intercostal space) can be used to better detect atrial activity in relation to that of the ventricles.
An esophageal lead can be inserted to a part of the esophagus where the distance to the posterior wall of the left atrium is only approximately 5–6 mm (remaining constant in people of different age and weight). An esophageal lead avails for a more accurate differentiation between certain cardiac arrhythmias, particularly atrial flutter, AV nodal reentrant tachycardia and orthodromic atrioventricular reentrant tachycardia. It can also evaluate the risk in people with Wolff-Parkinson-White syndrome, as well as terminate supraventricular tachycardia caused by re-entry.
An intracardiac electrogram (ICEG) is essentially an ECG with some added intracardiac leads (that is, inside the heart). The standard ECG leads (external leads) are I, II, III, aVL, V1, and V6. Two to four intracardiac leads are added via cardiac catheterization. The word "electrogram" (EGM) without further specification usually means an intracardiac electrogram.
Lead locations on an ECG report
A standard 12-lead ECG report (an electrocardiograph) shows a 2.5 second tracing of each of the twelve leads. The tracings are most commonly arranged in a grid of four columns and three rows. The first column is the limb leads (I, II, and III), the second column is the augmented limb leads (aVR, aVL, and aVF), and the last two columns are the precordial leads (V1 to V6).
Additionally, a rhythm strip may be included as a fourth or fifth row.
The timing across the page is continuous and notes tracings of the 12 leads for the same time period.
In other words, if the output were traced by needles on paper, each row would switch which leads as the paper is pulled under the needle.
For example, the top row would first trace lead I, then switch to lead aVR, then switch to V1, and then switch to V4, and so none of these four tracings of the leads are from the same time period as they are traced in sequence through time.
Contiguity of leads
Each of the 12 ECG leads records the electrical activity of the heart from a different angle, and therefore align with different anatomical areas of the heart. Two leads that look at neighboring anatomical areas are said to be contiguous.
In addition, any two precordial leads next to one another are considered to be contiguous. For example, though V4 is an anterior lead and V5 is a lateral lead, they are contiguous because they are next to one another.
Electrophysiology
The study of the conduction system of the heart is called cardiac electrophysiology (EP). An EP study is performed via a right-sided cardiac catheterization: a wire with an electrode at its tip is inserted into the right heart chambers from a peripheral vein, and placed in various positions in close proximity to the conduction system so that the electrical activity of that system can be recorded.
Standard catheter positions for an EP study include "high right atrium" or hRA near the sinus node, a "His" across the septal wall of the tricuspid valve to measure bundle of His, a "coronary sinus" into the coronary sinus, and a "right ventricle" in the apex of the right ventricle.
Interpretation
Interpretation of the ECG is fundamentally about understanding the electrical conduction system of the heart.
Normal conduction starts and propagates in a predictable pattern, and deviation from this pattern can be a normal variation or be pathological.
An ECG does not equate with mechanical pumping activity of the heart; for example, pulseless electrical activity produces an ECG that should pump blood but no pulses are felt (and constitutes a medical emergency and CPR should be performed).
Ventricular fibrillation produces an ECG but is too dysfunctional to produce a life-sustaining cardiac output. Certain rhythms are known to have good cardiac output and some are known to have bad cardiac output.
Ultimately, an echocardiogram or other anatomical imaging modality is useful in assessing the mechanical function of the heart.
Like all medical tests, what constitutes "normal" is based on population studies. The heartrate range of between 60 and 100 beats per minute (bpm) is considered normal since data shows this to be the usual resting heart rate.
Theory
Interpretation of the ECG is ultimately that of pattern recognition.
In order to understand the patterns found, it is helpful to understand the theory of what ECGs represent.
The theory is rooted in electromagnetics and boils down to the four following points:
depolarization of the heart toward the positive electrode produces a positive deflection
depolarization of the heart away from the positive electrode produces a negative deflection
repolarization of the heart toward the positive electrode produces a negative deflection
repolarization of the heart away from the positive electrode produces a positive deflection
Thus, the overall direction of depolarization and repolarization produces positive or negative deflection on each lead's trace.
For example, depolarizing from right to left would produce a positive deflection in lead I because the two vectors point in the same direction.
In contrast, that same depolarization would produce minimal deflection in V1 and V2 because the vectors are perpendicular, and this phenomenon is called isoelectric.
Normal rhythm produces four entities – a P wave, a QRS complex, a T wave, and a U wave – that each have a fairly unique pattern.
The P wave represents atrial depolarization.
The QRS complex represents ventricular depolarization.
The T wave represents ventricular repolarization.
The U wave represents papillary muscle repolarization.
Changes in the structure of the heart and its surroundings (including blood composition) change the patterns of these four entities.
The U wave is not typically seen and its absence is generally ignored. Atrial repolarization is typically hidden in the much more prominent QRS complex and normally cannot be seen without additional, specialized electrodes.
Background grid
ECGs are normally printed on a grid.
The horizontal axis represents time and the vertical axis represents voltage.
The standard values on this grid are shown in the adjacent image at 25mm/sec:
A small box is 1 mm × 1 mm and represents 0.1 mV × 0.04 seconds.
A large box is 5 mm × 5 mm and represents 0.5 mV × 0.20 seconds.
The "large" box is represented by a heavier line weight than the small boxes.
The standard printing speed in the United States is 25 mm per sec (5 big boxes per second), but in other countries it can be 50 mm per sec.
Faster speeds such as 100 and 200 mm per sec are used during electrophysiology studies.
Not all aspects of an ECG rely on precise recordings or having a known scaling of amplitude or time.
For example, determining if the tracing is a sinus rhythm only requires feature recognition and matching, and not measurement of amplitudes or times (i.e., the scale of the grids are irrelevant).
An example to the contrary, the voltage requirements of left ventricular hypertrophy require knowing the grid scale.
Rate and rhythm
In a normal heart, the heart rate is the rate at which the sinoatrial node depolarizes since it is the source of depolarization of the heart.
Heart rate, like other vital signs such as blood pressure and respiratory rate, change with age.
In adults, a normal heart rate is between 60 and 100 bpm (normocardic), whereas it is higher in children.
A heart rate below normal is called "bradycardia" (<60 in adults) and above normal is called "tachycardia" (>100 in adults).
A complication of this is when the atria and ventricles are not in synchrony and the "heart rate" must be specified as atrial or ventricular (e.g., the ventricular rate in ventricular fibrillation is 300–600 bpm, whereas the atrial rate can be normal [60–100] or faster [100–150]).
In normal resting hearts, the physiologic rhythm of the heart is normal sinus rhythm (NSR).
Normal sinus rhythm produces the prototypical pattern of P wave, QRS complex, and T wave.
Generally, deviation from normal sinus rhythm is considered a cardiac arrhythmia.
Thus, the first question in interpreting an ECG is whether or not there is a sinus rhythm.
A criterion for sinus rhythm is that P waves and QRS complexes appear 1-to-1, thus implying that the P wave causes the QRS complex.
Once sinus rhythm is established, or not, the second question is the rate.
For a sinus rhythm, this is either the rate of P waves or QRS complexes since they are 1-to-1.
If the rate is too fast, then it is sinus tachycardia, and if it is too slow, then it is sinus bradycardia.
If it is not a sinus rhythm, then determining the rhythm is necessary before proceeding with further interpretation.
Some arrhythmias with characteristic findings:
Absent P waves with "irregularly irregular" QRS complexes are the hallmark of atrial fibrillation.
A "saw tooth" pattern with QRS complexes is the hallmark of atrial flutter.
A sine wave pattern is the hallmark of ventricular flutter.
Absent P waves with wide QRS complexes and a fast heart rate are ventricular tachycardia.
Determination of rate and rhythm is necessary in order to make sense of further interpretation.
Axis
The heart has several axes, but the most common by far is the axis of the QRS complex (references to "the axis" imply the QRS axis).
Each axis can be computationally determined to result in a number representing degrees of deviation from zero, or it can be categorized into a few types.
The QRS axis is the general direction of the ventricular depolarization wavefront (or mean electrical vector) in the frontal plane.
It is often sufficient to classify the axis as one of three types: normal, left deviated, or right deviated.
Population data shows that a normal QRS axis is from −30° to 105°, with 0° being along lead I and positive being inferior and negative being superior (best understood graphically as the hexaxial reference system).
Beyond +105° is right axis deviation and beyond −30° is left axis deviation (the third quadrant of −90° to −180° is very rare and is an indeterminate axis).
A shortcut for determining if the QRS axis is normal is if the QRS complex is mostly positive in lead I and lead II (or lead I and aVF if +90° is the upper limit of normal).
The normal QRS axis is generally down and to the left, following the anatomical orientation of the heart within the chest. An abnormal axis suggests a change in the physical shape and orientation of the heart or a defect in its conduction system that causes the ventricles to depolarize in an abnormal way.
The extent of a normal axis can be +90° or 105° depending on the source.
Amplitudes and intervals
All of the waves on an ECG tracing and the intervals between them have a predictable time duration, a range of acceptable amplitudes (voltages), and a typical morphology. Any deviation from the normal tracing is potentially pathological and therefore of clinical significance.
For ease of measuring the amplitudes and intervals, an ECG is printed on graph paper at a standard scale: each 1 mm (one small box on the standard 25mm/s ECG paper) represents 40 milliseconds of time on the x-axis, and 0.1 millivolts on the y-axis.
Time-Frequency Analysis in ECG Signal Processing
In electrocardiogram (ECG) signal processing, Time-Frequency Analysis (TFA) is an important technique used to reveal how the frequency characteristics of ECG signals change over time, especially in non-stationary signals such as arrhythmias or transient cardiac events.
Common Methods for Time-Frequency Analysis
Steps for Time-Frequency Analysis
Step1: Preprocessing
Signal Denoising: Use wavelet denoising, band-pass filtering (0.5–50 Hz), or Principal Component Analysis (PCA) to remove electromyographic (EMG) noise.
Signal Segmentation: Segment the signal based on heartbeat cycles (e.g., R-wave detection).
Step2: Select an Appropriate TFA Method
Choose methods such as STFT, WT, or HHT based on the application requirements.
Step3: Compute the Time-Frequency Spectrum
Calculate the time-frequency distribution using the selected method to generate a time-frequency representation.
Step4: Feature Extraction
Extract power features from specific frequency bands, such as low-frequency (LF: 0.04–0.15 Hz) and high-frequency (HF: 0.15–0.4 Hz) components.
Step5: Pattern Recognition or Diagnosis
Apply machine learning or deep learning models to detect or classify cardiac events based on the time-frequency features.
Application Scenarios
Heart Rate Variability Analysis (HRV):
Time-frequency analysis helps to separate sympathetic and parasympathetic nervous system activity.
Atrial Fibrillation Detection:
Analyze the time-frequency characteristics of atrial activity.
Ventricular Fibrillation Analysis:
Detect time-frequency changes in high-frequency abnormal components.
Limb leads and electrical conduction through the heart
The animation shown to the right illustrates how the path of electrical conduction gives rise to the ECG waves in the limb leads.
What is green zone ?
Recall that a positive current (as created by depolarization of cardiac cells) traveling towards the positive electrode and away from the negative electrode creates a positive deflection on the ECG. Likewise, a positive current traveling away from the positive electrode and towards the negative electrode creates a negative deflection on the ECG. The red arrow represents the overall direction of travel of the depolarization. The magnitude of the red arrow is proportional to the amount of tissue being depolarized at that instance. The red arrow is simultaneously shown on the axis of each of the 3 limb leads. Both the direction and the magnitude of the red arrow's projection onto the axis of each limb lead is shown with blue arrows. Then, the direction and magnitude of the blue arrows are what theoretically determine the deflections on the ECG. For example, as a blue arrow on the axis for Lead I moves from the negative electrode, to the right, towards the positive electrode, the ECG line rises, creating an upward wave. As the blue arrow on the axis for Lead I moves to the left, a downward wave is created. The greater the magnitude of the blue arrow, the greater the deflection on the ECG for that particular limb lead.
Frames 1–3 depict the depolarization being generated in and spreading through the sinoatrial node. The SA node is too small for its depolarization to be detected on most ECGs. Frames 4–10 depict the depolarization traveling through the atria, towards the atrioventricular node. During frame 7, the depolarization is traveling through the largest amount of tissue in the atria, which creates the highest point in the P wave. Frames 11–12 depict the depolarization traveling through the AV node. Like the SA node, the AV node is too small for the depolarization of its tissue to be detected on most ECGs. This creates the flat PR segment.
Frame 13 depicts an interesting phenomenon in an over-simplified fashion. It depicts the depolarization as it starts to travel down the interventricular septum, through the bundle of His and bundle branches. After the Bundle of His, the conduction system splits into the left bundle branch and the right bundle branch. Both branches conduct action potentials at about 1 m/s. However, the action potential starts traveling down the left bundle branch about 5 milliseconds before it starts traveling down the right bundle branch, as depicted by frame 13. This causes the depolarization of the interventricular septum tissue to spread from left to right, as depicted by the red arrow in frame 14. In some cases, this gives rise to a negative deflection after the PR interval, creating a Q wave such as the one seen in lead I in the animation to the right. Depending on the mean electrical axis of the heart, this phenomenon can result in a Q wave in lead II as well.
Following depolarization of the interventricular septum, the depolarization travels towards the apex of the heart. This is depicted by frames 15–17 and results in a positive deflection on all three limb leads, which creates the R wave. Frames 18–21 then depict the depolarization as it travels throughout both ventricles from the apex of the heart, following the action potential in the Purkinje fibers. This phenomenon creates a negative deflection in all three limb leads, forming the S wave on the ECG. Repolarization of the atria occurs at the same time as the generation of the QRS complex, but it is not detected by the ECG since the tissue mass of the ventricles is so much larger than that of the atria. Ventricular contraction occurs between ventricular depolarization and repolarization. During this time, there is no movement of charge, so no deflection is created on the ECG. This results in the flat ST segment after the S wave.
Frames 24–28 in the animation depict repolarization of the ventricles. The epicardium is the first layer of the ventricles to repolarize, followed by the myocardium. The endocardium is the last layer to repolarize. The plateau phase of depolarization has been shown to last longer in endocardial cells than in epicardial cells. This causes repolarization to start from the apex of the heart and move upwards. Since repolarization is the spread of negative current as membrane potentials decrease back down to the resting membrane potential, the red arrow in the animation is pointing in the direction opposite of the repolarization. This therefore creates a positive deflection in the ECG, and creates the T wave.
Ischemia and infarction
Ischemia or non-ST elevation myocardial infarctions (non-STEMIs) may manifest as ST depression or inversion of T waves. It may also affect the high frequency band of the QRS.
ST elevation myocardial infarctions (STEMIs) have different characteristic ECG findings based on the amount of time elapsed since the MI first occurred. The earliest sign is hyperacute T waves, peaked T waves due to local hyperkalemia in ischemic myocardium. This then progresses over a period of minutes to elevations of the ST segment by at least 1 mm. Over a period of hours, a pathologic Q wave may appear and the T wave will invert. Over a period of days the ST elevation will resolve. Pathologic Q waves generally will remain permanently.
The coronary artery that has been occluded can be identified in an STEMI based on the location of ST elevation. The left anterior descending (LAD) artery supplies the anterior wall of the heart, and therefore causes ST elevations in anterior leads (V1 and V2). The LCx supplies the lateral aspect of the heart and therefore causes ST elevations in lateral leads (I, aVL and V6). The right coronary artery (RCA) usually supplies the inferior aspect of the heart, and therefore causes ST elevations in inferior leads (II, III and aVF).
Artifacts
An ECG tracing is affected by patient motion. Some rhythmic motions (such as shivering or tremors) can create the illusion of cardiac arrhythmia. Artifacts are distorted signals caused by a secondary internal or external sources, such as muscle movement or interference from an electrical device.
Distortion poses significant challenges to healthcare providers, who employ various techniques and strategies to safely recognize these false signals. Accurately separating the ECG artifact from the true ECG signal can have a significant impact on patient outcomes and legal liabilities.
Improper lead placement (for example, reversing two of the limb leads) has been estimated to occur in 0.4% to 4% of all ECG recordings, and has resulted in improper diagnosis and treatment including unnecessary use of thrombolytic therapy.
A Method for Interpretation
Whitbread, consultant nurse and paramedic, suggests ten rules of the normal ECG, deviation from which is likely to indicate pathology. These have been added to, creating the 15 rules for 12-lead (and 15- or 18-lead) interpretation.
Rule 1: All waves in aVR are negative.
Rule 2: The ST segment (J point) starts on the isoelectric line (except in V1 & V2 where it may be elevated by not greater than 1 mm).
Rule 3: The PR interval should be 0.12–0.2 seconds long.
Rule 4: The QRS complex should not exceed 0.11–0.12 seconds.
Rule 5: The QRS and T waves tend to have the same general direction in the limb leads.
Rule 6: The R wave in the precordial (chest) leads grows from V1 to at least V4 where it may or may not decline again.
Rule 7: The QRS is mainly upright in I and II.
Rule 8: The P wave is upright in I II and V2 to V6.
Rule 9: There is no Q wave or only a small q (<0.04 seconds in width) in I, II and V2 to V6.
Rule 10: The T wave is upright in I II and V2 to V6. The end of the T wave should not drop below the isoelectric baseline.
Rule 11: Does the deepest S wave in V1 plus the tallest R wave in V5 or V6 equal >35 mm?
Rule 12: Is there an Epsilon wave?
Rule 13: Is there an J wave?
Rule 14: Is there a Delta wave?
Rule 15: Are there any patterns representing an occlusive myocardial infarction (OMI)?
Diagnosis
Numerous diagnoses and findings can be made based upon electrocardiography, and many are discussed above. Overall, the diagnoses are made based on the patterns. For example, an "irregularly irregular" QRS complex without P waves is the hallmark of atrial fibrillation; however, other findings can be present as well, such as a bundle branch block that alters the shape of the QRS complexes. ECGs can be interpreted in isolation but should be applied – like all diagnostic tests – in the context of the patient. For example, an observation of peaked T waves is not sufficient to diagnose hyperkalemia; such a diagnosis should be verified by measuring the blood potassium level. Conversely, a discovery of hyperkalemia should be followed by an ECG for manifestations such as peaked T waves, widened QRS complexes, and loss of P waves. The following is an organized list of possible ECG-based diagnoses.
Rhythm disturbances or arrhythmias:
Atrial fibrillation and atrial flutter without rapid ventricular response
Premature atrial contraction (PACs) and premature ventricular contraction (PVCs)
Sinus arrhythmia
Sinus bradycardia and sinus tachycardia
Sinus pause and sinoatrial arrest
Sinus node dysfunction and bradycardia-tachycardia syndrome
Supraventricular tachycardia
Atrial fibrillation with rapid ventricular response
Atrial flutter with rapid ventricular response
AV nodal reentrant tachycardia
Atrioventricular reentrant tachycardia
Junctional ectopic tachycardia
Atrial tachycardia
Ectopic atrial tachycardia (unicentric)
Multifocal atrial tachycardia
Paroxysmal atrial tachycardia
Sinoatrial nodal reentrant tachycardia
Wide complex tachycardia
Ventricular flutter
Ventricular fibrillation
Ventricular tachycardia (monomorphic ventricular tachycardia)
Torsades de pointes (polymorphic ventricular tachycardia)
Pre-excitation syndrome
Lown–Ganong–Levine syndrome
Wolff–Parkinson–White syndrome
J wave (Osborn wave)
Heart block and conduction problems:
Sinoatrial block: first, second, and third-degree
AV node
First-degree AV block
Second-degree AV block (Mobitz [Wenckebach] I and II)
Third-degree AV block or complete AV block
Right bundle
Incomplete right bundle branch block (IRBBB)
Complete right bundle branch block (RBBB)
Left bundle
Incomplete left bundle branch block (ILBBB)
Complete left bundle branch block (LBBB)
Left anterior fascicular block (LAFB)
Left posterior fascicular block (LPFB)
Bifascicular block (LAFB plus LPFB)
Trifascicular block (LAFP plus FPFB plus RBBB)
QT syndromes
Brugada syndrome
Short QT syndrome
Long QT syndromes, genetic and drug-induced
Right and left atrial abnormality
Electrolytes disturbances and intoxication:
Digitalis intoxication
Calcium: hypocalcemia and hypercalcemia
Potassium: hypokalemia and hyperkalemia
Serotonin toxicity
Ischemia and infarction:
Wellens' syndrome (LAD occlusion)
de Winter T waves (LAD occlusion)
ST elevation and ST depression
High frequency QRS changes
Myocardial infarction (heart attack)
Non-Q wave myocardial infarction
NSTEMI
STEMI
Sgarbossa's criteria for ischemia with a LBBB
Structural:
Acute pericarditis
Right and left ventricular hypertrophy
Right ventricular strain or S1Q3T3 (can be seen in pulmonary embolism)
Other phenomena:
Cardiac aberrancy
Ashman phenomenon
Concealed conduction
Electrical alternans
History
In 1872, Alexander Muirhead is reported to have attached wires to the wrist of a patient with fever to obtain an electronic record of their heartbeat.
In 1882, John Burdon-Sanderson working with frogs, was the first to appreciate that the interval between variations in potential was not electrically quiescent and coined the term "isoelectric interval" for this period.
In 1887, Augustus Waller invented an ECG machine consisting of a Lippmann capillary electrometer fixed to a projector. The trace from the heartbeat was projected onto a photographic plate that was itself fixed to a toy train. This allowed a heartbeat to be recorded in real time.
In 1895, Willem Einthoven assigned the letters P, Q, R, S, and T to the deflections in the theoretical waveform he created using equations which corrected the actual waveform obtained by the capillary electrometer to compensate for the imprecision of that instrument. Using letters different from A, B, C, and D (the letters used for the capillary electrometer's waveform) facilitated comparison when the uncorrected and corrected lines were drawn on the same graph. Einthoven probably chose the initial letter P to follow the example set by Descartes in geometry. When a more precise waveform was obtained using the string galvanometer, which matched the corrected capillary electrometer waveform, he continued to use the letters P, Q, R, S, and T, and these letters are still in use today. Einthoven also described the electrocardiographic features of a number of cardiovascular disorders.
In 1897, the string galvanometer was invented by the French engineer Clément Ader.
In 1901, Einthoven, working in Leiden, the Netherlands, used the string galvanometer: the first practical ECG. This device was much more sensitive than the capillary electrometer Waller used.
In 1924, Einthoven was awarded the Nobel Prize in Medicine for his pioneering work in developing the ECG.
By 1927, General Electric had developed a portable apparatus that could produce electrocardiograms without the use of the string galvanometer. This device instead combined amplifier tubes similar to those used in a radio with an internal lamp and a moving mirror that directed the tracing of the electric pulses onto film.
In 1937, Taro Takemi invented a new portable electrocardiograph machine.
In 1942, Emanuel Goldberger increases the voltage of Wilson's unipolar leads by 50% and creates the augmented limb leads aVR, aVL and aVF. When added to Einthoven's three limb leads and the six chest leads we arrive at the 12-lead electrocardiogram that is used today.
In the late 1940s, Rune Elmqvist invented an inkjet printer involving thin jets of ink deflected by electrical potentials from the heart, with good frequency response and direct recording of ECG on paper. The device, called the Mingograf, was sold by Siemens Elema until the 1990s.
Etymology
The word is derived from the Greek electro, meaning related to electrical activity; kardia, meaning heart; and graph, meaning "to write".
| Technology | Diagnostic technologies | null |
76992 | https://en.wikipedia.org/wiki/Tartaric%20acid | Tartaric acid | Tartaric acid is a white, crystalline organic acid that occurs naturally in many fruits, most notably in grapes but also in tamarinds, bananas, avocados, and citrus. Its salt, potassium bitartrate, commonly known as cream of tartar, develops naturally in the process of fermentation. Potassium bitartrate is commonly mixed with sodium bicarbonate and is sold as baking powder used as a leavening agent in food preparation. The acid itself is added to foods as an antioxidant E334 and to impart its distinctive sour taste. Naturally occurring tartaric acid is a useful raw material in organic synthesis. Tartaric acid, an alpha-hydroxy-carboxylic acid, is diprotic and aldaric in acid characteristics and is a dihydroxyl derivative of succinic acid.
History
Tartaric acid has been known to winemakers for centuries. However, the chemical process for extraction was developed in 1769 by the Swedish chemist Carl Wilhelm Scheele.
Tartaric acid played an important role in the discovery of chemical chirality. This property of tartaric acid was first observed in 1832 by Jean Baptiste Biot, who observed its ability to rotate polarized light. Louis Pasteur continued this research in 1847 by investigating the shapes of sodium ammonium tartrate crystals, which he found to be chiral. By manually sorting the differently shaped crystals, Pasteur was the first to produce a pure sample of levotartaric acid.
Stereochemistry
Naturally occurring form of the acid is dextro tartaric acid or L-(+)-tartaric acid (obsolete name d-tartaric acid). Because it is available naturally, it is cheaper than its enantiomer and the meso isomer. The dextro and levo prefixes are archaic terms. Modern textbooks refer to the natural form as (2R,3R)-tartaric acid (L-(+)-tartaric acid), and its enantiomer as (2S,3S)-tartaric acid (D-(-)-tartaric acid). The meso diastereomer is referred to as (2R,3S)-tartaric acid or (2S,3R)-tartaric acid.
Dextro and levo form monoclinic sphenoidal crystals and orthorhombic crystals.
Racemic tartaric acid forms monoclinic and triclinic crystals (space group P).
Anhydrous meso tartaric acid form two anhydrous polymorphs: triclinic and orthorhombic.
Monohydrated meso tartaric acid crystallizes as monoclinic and triclinic polymorphys depending on the temperature at which crystallization from aqueous solution occurs.
Tartaric acid in Fehling's solution binds to copper(II) ions, preventing the formation of insoluble hydroxide salts.
Production
L-(+)-Tartaric acid
The L-(+)-tartaric acid isomer of tartaric acid is industrially produced in the largest amounts. It is obtained from lees, a solid byproduct of fermentations. The former byproducts mostly consist of potassium bitartrate (). This potassium salt is converted to calcium tartrate () upon treatment with calcium hydroxide ():
In practice, higher yields of calcium tartrate are obtained with the addition of calcium sulfate. Calcium tartrate is then converted to tartaric acid by treating the salt with aqueous sulfuric acid:
Racemic tartaric acid
Racemic tartaric acid can be prepared in a multistep reaction from maleic acid. In the first step, the maleic acid is epoxidized by hydrogen peroxide using as a catalyst.
In the next step, the epoxide is hydrolyzed.
meso-Tartaric acid
A mixture of racemic acid and meso-tartaric acid is formed when dextro-Tartaric acid is heated in water at 165 °C for about 2 days. meso-Tartaric acid can also be prepared from dibromosuccinic acid using silver hydroxide:
meso-Tartaric acid can be separated from residual racemic acid by crystallization, the racemate being less soluble.
Reactivity
L-(+)-tartaric acid, can participate in several reactions. As shown the reaction scheme below, dihydroxymaleic acid is produced upon treatment of L-(+)-tartaric acid with hydrogen peroxide in the presence of a ferrous salt.
HO2CCH(OH)CH(OH)CO2H + H2O2 → HO2CC(OH)C(OH)CO2H + 2 H2O
Dihydroxymaleic acid can then be oxidized to tartronic acid with nitric acid.
Derivatives
Important derivatives of tartaric acid include:
Sodium ammonium tartrate, the first material separated into its enantiomers
cream of tartar (potassium bitartrate), used in cooking
Rochelle salt (potassium sodium tartrate), which has unusual piezoelectric properties
tartar emetic (antimony potassium tartrate), a resolving agent. Diisopropyl tartrate is used as a co-catalyst in asymmetric synthesis.
Tartaric acid is a muscle toxin, which works by inhibiting the production of malic acid, and in high doses causes paralysis and death. The median lethal dose (LD50) is about 7.5 grams/kg for a human, 5.3 grams/kg for rabbits, and 4.4 grams/kg for mice. Given this figure, it would take over to kill a person weighing with 50% probability, so it may be safely included in many foods, especially sour-tasting sweets. As a food additive, tartaric acid is used as an antioxidant with E number E334; tartrates are other additives serving as antioxidants or emulsifiers.
When cream of tartar is added to water, a suspension results which serves to clean copper coins very well, as the tartrate solution can dissolve the layer of copper(II) oxide present on the surface of the coin. The resulting copper(II)-tartrate complex is easily soluble in water.
Tartaric acid in wine
Tartaric acid may be most immediately recognizable to wine drinkers as the source of "wine diamonds", the small potassium bitartrate crystals that sometimes form spontaneously on the cork or bottom of the bottle. These "tartrates" are harmless, despite sometimes being mistaken for broken glass, and are prevented in many wines through cold stabilization (which is not always preferred since it can change the wine's profile). The tartrates remaining on the inside of aging barrels were at one time a major industrial source of potassium bitartrate.
Tartaric acid plays an important role chemically, lowering the pH of fermenting "must" to a level where many undesirable spoilage bacteria cannot live, and acting as a preservative after fermentation. In the mouth, tartaric acid provides some of the tartness in the wine, although citric and malic acids also play a role.
Tartaric acid in fruits
Grapes and tamarinds have the highest levels of tartaric acid concentration. Other fruits with tartaric acid are bananas, avocados, prickly pear fruit, apples, cherries, papayas, peaches, pears, pineapples, strawberries, mangoes and citrus fruits.
Trace amounts of tartaric acid have been found in cranberries and other berries.
Tartaric acid is also present in the leaves and pods of Pelargonium plants and beans.
Applications
Tartaric acid and its derivatives have a plethora of uses in the field of pharmaceuticals. For example, it has been used in the production of effervescent salts, in combination with citric acid, to improve the taste of oral medications. The potassium antimonyl derivative of the acid known as tartar emetic is included, in small doses, in cough syrup as an expectorant.
Tartaric acid also has several applications for industrial use. The acid has been observed to chelate metal ions such as calcium and magnesium. Therefore, the acid has served in the farming and metal industries as a chelating agent for complexing micronutrients in soil fertilizer and for cleaning metal surfaces consisting of aluminium, copper, iron, and alloys of these metals, respectively.
Toxicity in canines
While tartaric acid is well-tolerated by humans and lab animals, an April 2021 letter to the editor of JAVMA hypothesized that the tartaric acid in grapes could be the cause of grape and raisin toxicity in dogs. Other studies have observed tartaric acid toxicity in kidney cells of dogs, but not in human kidney cells.
In dogs, the tartaric acid of tamarind causes acute kidney injury, which can often be fatal.
A review identified a relationship between grape ingestion and illness, though the specific type or quantity of grapes that cause toxicity remains unclear. Grape ingestion commonly leads to gastrointestinal and/or renal issues, with treatment depending on the symptoms; outcomes can vary.
| Physical sciences | Specific acids | Chemistry |
77060 | https://en.wikipedia.org/wiki/Biological%20hazard | Biological hazard | A biological hazard, or biohazard, is a biological substance that poses a threat (or is a hazard) to the health of living organisms, primarily humans. This could include a sample of a microorganism, virus or toxin that can adversely affect human health. A biohazard could also be a substance harmful to other living beings.
The term and its associated symbol are generally used as a warning, so that those potentially exposed to the substances will know to take precautions. The biohazard symbol was developed in 1966 by Charles Baldwin, an environmental-health engineer working for the Dow Chemical Company on their containment products. It is used in the labeling of biological materials that carry a significant health risk, including viral samples and used hypodermic needles. In Unicode, the biohazard symbol is U+2623 (☣).
ANSI Z535/OSHA/ISO regulation
Biohazardous safety issues are identified with specified labels, signs and paragraphs established by the American National Standards Institute (ANSI). Today, ANSI Z535 standards for biohazards are used worldwide and should always be used appropriately within ANSI Z535 Hazardous Communications (HazCom) signage, labeling and paragraphs. The goal is to help workers rapidly identify the severity of a biohazard from a distance and through colour and design standardization.
Biological hazard symbol design:
A red on white or white-coloured background is used behind a black biohazard symbol when integrated with a DANGER sign, label or paragraph.
An orange on black or white-coloured background is used behind a black biohazard symbol when integrated with a WARNING sign, label or paragraph.
A yellow on black or white-coloured background is used behind a black biohazard symbol when integrated with a CAUTION sign, label or paragraph.
A green on white or white-coloured background is used behind a black biohazard symbol when integrated with a NOTICE sign, label or paragraph.
DANGER is used to identify a biohazard that will cause death. WARNING is used to identify a biohazard that may cause death. CAUTION is used to identify a biohazard that will cause injury, but not death. NOTICE is used to identify a non-injury biohazard message (e.g. hygiene, cleanup or general lab policies).
OSHA requires the use of proper ANSI HazCom where applicable in American workplaces. States and local governments also use these standards as codes and laws within their own jurisdictions. Proper use of ANSI Z535 signs, labels and paragraphs are written into many of OSHA's standards for HazCom and crafted to integrate with ISO symbols.
Reference ANSI Z535 for a complete description on how to use DANGER, WARNING, CAUTION and NOTICE signs, labels or paragraphs.
UN/ISO Classification
Biohazardous agents are classified for transportation by UN number:
Category A, UN 2814 – Infectious substance, affecting humans: An infectious substance in a form capable of causing permanent disability or life-threatening or fatal disease in otherwise healthy humans or animals when exposure to it occurs.
Category A, UN 2900 – Infectious substance, affecting animals (only): An infectious substance that is not in a form generally capable of causing permanent disability or life-threatening or fatal disease in otherwise healthy humans and animals when exposure to themselves occurs.
Category B, UN 3373 – Biological substance transported for diagnostic or investigative purposes.
Regulated Medical Waste, UN 3291 – Waste or reusable material derived from medical treatment of an animal or human, or from biomedical research, which includes the production and testing.
Levels of biohazard
The United States Centers for Disease Control and Prevention (CDC) categorizes various diseases in levels of biohazard, Level 1 being minimum risk and Level 4 being extreme risk. Laboratories and other facilities are categorized as BSL (Biosafety Level) 1–4 or as P1 through P4 for short (Pathogen or Protection Level).
Biohazard Level 1: Bacteria and viruses including Bacillus subtilis, canine hepatitis, Escherichia coli, and varicella (chickenpox), as well as some cell cultures and non-infectious bacteria. At this level precautions against the biohazardous materials in question are minimal, most likely involving gloves and some sort of facial protection.
Biohazard Level 2: Bacteria and viruses that cause only mild disease to humans, or are difficult to contract via aerosol in a lab setting, such as hepatitis A, B, and C, some influenza A strains, Human respiratory syncytial virus, Lyme disease, salmonella, mumps, measles, scrapie, dengue fever, and HIV. Routine diagnostic work with clinical specimens can be done safely at Biosafety Level 2, using Biosafety Level 2 practices and procedures. Research work (including co-cultivation, virus replication studies, or manipulations involving concentrated virus) can be done in a BSL-2 (P2) facility, using BSL-3 practices and procedures.
Biohazard Level 3: Bacteria and viruses that can cause severe to fatal disease in humans, but for which vaccines or other treatments exist, such as anthrax, West Nile virus, Venezuelan equine encephalitis, SARS coronavirus, MERS coronavirus, SARS-CoV-2, Influenza A H5N1, hantaviruses, Cholera, tuberculosis, typhus, Rift Valley fever, Rocky Mountain spotted fever, yellow fever, and malaria.
Biohazard Level 4: Viruses that cause severe to fatal disease in humans, and for which vaccines or other treatments are not available, such as Bolivian hemorrhagic fever, Marburg virus, Ebola virus, Lassa fever virus, Crimean–Congo hemorrhagic fever, and other hemorrhagic diseases, as well as Nipah virus. Variola virus (smallpox) is an agent that is worked with at BSL-4 despite the existence of a vaccine, as it has been eradicated and thus the general population is no longer routinely vaccinated. When dealing with biological hazards at this level, the use of a positive pressure personnel suit with a segregated air supply is mandatory. The entrance and exit of a Level Four biolab will contain multiple showers, a vacuum room, an ultraviolet light room, autonomous detection system, and other safety precautions designed to destroy all traces of the biohazard. Multiple airlocks are employed and are electronically secured to prevent doors from both opening at the same time. All air and water service going to and coming from a Biosafety Level 4 (P4) lab will undergo similar decontamination procedures to eliminate the possibility of an accidental release. Currently there are no bacteria classified at this level.
| Technology | Biotechnology | null |
77178 | https://en.wikipedia.org/wiki/Spaceflight | Spaceflight | Spaceflight (or space flight) is an application of astronautics to fly objects, usually spacecraft, into or through outer space, either with or without humans on board. Most spaceflight is uncrewed and conducted mainly with spacecraft such as satellites in orbit around Earth, but also includes space probes for flights beyond Earth orbit. Such spaceflights operate either by telerobotic or autonomous control. The first spaceflights began in the 1950s with the launches of the Soviet Sputnik satellites and American Explorer and Vanguard missions. Human spaceflight programs include the Soyuz, Shenzhou, the past Apollo Moon landing and the Space Shuttle programs. Other current spaceflight are conducted to the International Space Station and to China's Tiangong Space Station.
Spaceflights include the launches of Earth observation and telecommunications satellites, interplanetary missions, the rendezvouses and dockings with space stations, and crewed spaceflights on scientific or tourist missions.
Spaceflight can be achieved conventionally via multistage rockets, which provide the thrust to overcome the force of gravity and propel spacecraft onto suborbital trajectories. If the mission is orbital, the spacecraft usually separates the first stage and ignites the second stage, which propels the spacecraft to high enough speeds that it reaches orbit. Once in orbit, spacecraft are at high enough speeds that they fall around the Earth rather than fall back to the surface.
Most spacecraft, and all crewed spacecraft, are designed to deorbit themselves or, in the case of uncrewed spacecraft in high-energy orbits, to boost themselves into graveyard orbits. Used upper stages or failed spacecraft, however, often lack the ability to deorbit themselves. This becomes a major issue when large numbers of uncontrollable spacecraft exist in frequently used orbits, increasing the risk of debris colliding with functional satellites. This problem is exacerbated when large objects, often upper stages, break up in orbit or collide with other objects, creating often hundreds of small, hard to find pieces of debris. This problem of continuous collisions is known as Kessler syndrome.
Terminology
There are several terms that refer to a flight into or through outer space.
A space mission refers to a spaceflight intended to achieve an objective. Objectives for space missions may include space exploration, space research, and national firsts in spaceflight.
Space transport is the use of spacecraft to transport people or cargo into or through outer space. This may include human spaceflight and cargo spacecraft flight.
History
The first theoretical proposal of space travel using rockets was published by Scottish astronomer and mathematician William Leitch, in an 1861 essay "A Journey Through Space". More well-known is Konstantin Tsiolkovsky's work, "" (The Exploration of Cosmic Space by Means of Reaction Devices), published in 1903. In his work, Tsiolkovsky describes the fundamental rocket equation:
Where:
() is the change in the rocket's velocity
() is the exhaust velocity
() and () are the initial and final masses of the rocket
This equation, known as the Tsiolkovsky rocket equation, can be used to find the total , or potential change in velocity. This formula, which is still used by engineers, is a key concept of spaceflight.
Spaceflight became a practical possibility with the work of Robert H. Goddard's publication in 1919 of his paper A Method of Reaching Extreme Altitudes. His application of the de Laval nozzle to liquid-fuel rockets improved efficiency enough for interplanetary travel to become possible. After further research, Goddard attempted to secure an Army contract for a rocket-propelled weapon in the first World War but his plans were foiled by the November 11, 1918 armistice with Germany. After choosing to work with private financial support, he was the first to launch a liquid-fueled rocket on March 16, 1926.
During World War II, the first guided rocket, the V-2, was developed and employed as a weapon by Nazi Germany. During a test flight in June 1944, one such rocket reached space at an altitude of , becoming the first human-made object to reach space. At the end of World War II, most of the V-2 rocket team, including its head, Wernher von Braun, surrendered to the United States, and were expatriated to work on American missiles at what became the Army Ballistic Missile Agency, producing missiles such as Juno I and Atlas. The Soviet Union, in turn, captured several V2 production facilities and built several replicas, with 5 of their 11 rockets successfully reaching their targets. (This was relatively consistent with Nazi Germany's success rate.)
The Soviet Union developed intercontinental ballistic missiles to carry nuclear weapons as a counter measure to United States bomber planes in the 1950s. The Tsiolkovsky-influenced Sergey Korolev became the chief rocket designer, and derivatives of his R-7 Semyorka missiles were used to launch the world's first artificial Earth satellite, Sputnik 1, on October 4, 1957.
The U.S., after the launch of Sputnik and two embarrassing failures of Vanguard rockets, launched Explorer 1 on February 1, 1958. Three years later, the USSR launched Vostok 1, carrying cosmonaut Yuri Gagarin into orbit. The US responded with the suborbital launch of Alan Shepard on May 5, 1961, and the orbital launch of John Glenn on February 20, 1962. These events were followed by a pledge from U.S. President John F. Kennedy to go to the moon and the creation of the Gemini and Apollo programs. After successfully performing a rendezvous and docking and an EVA, the Gemini program ended just before the Apollo 1 tragedy. Following multiple uncrewed test flights of the Saturn 1B and the Saturn V, the U.S. launched the crewed Apollo 7 mission into low earth orbit. Shortly after its successful completion, the U.S. launched Apollo 8 (first mission to orbit the moon), Apollo 9 (first Apollo mission to launch with both the CSM and the LEM) and Apollo 10 (first mission to nearly land on the moon). These events culminated with the first crewed moon landing, Apollo 11, and six subsequent missions, five of which successfully landed on the moon.
Spaceflight has been widely employed by numerous government and commercial entities for placing satellites into orbit around Earth for a broad range of purposes. Certain government agencies have also sent uncrewed spacecraft exploring space beyond the Moon and developed continuous crewed human presence in space with a series of space stations, ranging from the Salyut program to the International Space Station.
Phases
Launch
Rockets are the only means currently capable of reaching orbit or beyond. Other non-rocket spacelaunch technologies have yet to be built, or remain short of orbital speeds. A rocket launch for a spaceflight usually starts from a spaceport (cosmodrome), which may be equipped with launch complexes and launch pads for vertical rocket launches and runways for takeoff and landing of carrier airplanes and winged spacecraft. Spaceports are situated well away from human habitation for noise and safety reasons. ICBMs have various special launching facilities.
A launch is often restricted to certain launch windows. These windows depend upon the position of celestial bodies and orbits relative to the launch site. The biggest influence is often the rotation of the Earth. Once launched, orbits are normally located within relatively constant flat planes at a fixed angle to the axis of the Earth, and the Earth rotates within this orbit.
A launch pad is a fixed structure designed to dispatch airborne vehicles. It generally consists of a launch tower and flame trench. It is surrounded by equipment used to erect, fuel, and maintain launch vehicles. Before launch, the rocket can weigh hundreds of tons. The Space Shuttle Columbia, on STS-1, weighed 2030 metric tons (4,480,000 lb) at takeoff.
Reaching space
The most commonly used definition of outer space is everything beyond the Kármán line, which is above the Earth's surface. (The United States defines outer space as everything beyond in altitude.)
Rocket engines remain the only currently practical means of reaching space, with planes and high-altitude balloons failing due to lack of atmosphere and alternatives such as space elevators not yet being built. Chemical propulsion, or the acceleration of gases at high velocities, is effective mainly because of its ability to sustain thrust even as the atmosphere thins.
Alternatives
Many ways to reach space other than rocket engines have been proposed. Ideas such as the space elevator, and momentum exchange tethers like rotovators or skyhooks require new materials much stronger than any currently known. Electromagnetic launchers such as launch loops might be feasible with current technology. Other ideas include rocket-assisted aircraft/spaceplanes such as Reaction Engines Skylon (currently in early stage development), scramjet powered spaceplanes, and RBCC powered spaceplanes. Gun launch has been proposed for cargo.
Leaving orbit
On some missions beyond LEO (Low Earth Orbit), spacecraft are inserted into parking orbits, or lower intermediary orbits. The parking orbit approach greatly simplified Apollo mission planning in several important ways. It acted as a "time buffer" and substantially widened the allowable launch windows. The parking orbit gave the crew and controllers time to thoroughly check out the spacecraft after the stresses of launch before committing it for a long journey to the Moon.
Robotic missions do not require an abort capability and require radiation minimalization only for delicate electronics, and because modern launchers routinely meet "instantaneous" launch windows, space probes to the Moon and other planets generally use direct injection to maximize performance by limiting the boil off of cryogenic propellants. Although some might coast briefly during the launch sequence, they do not complete one or more full parking orbits before the burn that injects them onto an Earth escape trajectory.
The escape velocity from a celestial body decreases as the distance from the body increases. However, it is more fuel-efficient for a craft to burn its fuel as close as possible to its periapsis (lowest point); see Oberth effect.
Astrodynamics
Astrodynamics is the study of spacecraft trajectories, particularly as they relate to gravitational and propulsion effects. Astrodynamics allows for a spacecraft to arrive at its destination at the correct time without excessive propellant use. An orbital maneuvering system may be needed to maintain or change orbits.
Non-rocket orbital propulsion methods include solar sails, magnetic sails, plasma-bubble magnetic systems, and using gravitational slingshot effects.
Transfer energy
The term "transfer energy" means the total amount of energy imparted by a rocket stage to its payload. This can be the energy imparted by a first stage of a launch vehicle to an upper stage plus payload, or by an upper stage or spacecraft kick motor to a spacecraft.
Reaching space station
In order to reach a space station, a spacecraft would have to arrive at the same orbit and approach to a very close distance (e.g. within visual contact). This is done by a set of orbital maneuvers called space rendezvous.
After rendezvousing with the space station, the space vehicle then docks or berths with the station. Docking refers to joining of two separate free-flying space vehicles, while berthing refers to mating operations where an inactive vehicle is placed into the mating interface of another space vehicle by using a robotic arm.
Reentry
Vehicles in orbit have large amounts of kinetic energy. This energy must be discarded if the vehicle is to land safely without vaporizing in the atmosphere. Typically this process requires special methods to protect against aerodynamic heating. The theory behind reentry was developed by Harry Julian Allen. Based on this theory, reentry vehicles present blunt shapes to the atmosphere for reentry. Blunt shapes mean that less than 1% of the kinetic energy ends up as heat reaching the vehicle, and the remainder heats the atmosphere.
Landing and recovery
The Mercury, Gemini, and Apollo capsules splashed down in the sea. These capsules were designed to land at relatively low speeds with the help of a parachute. Soviet/Russian capsules for Soyuz make use of a big parachute and braking rockets to touch down on land. Spaceplanes like the Space Shuttle land like a glider.
After a successful landing, the spacecraft, its occupants, and cargo can be recovered. In some cases, recovery has occurred before landing: while a spacecraft is still descending on its parachute, it can be snagged by a specially designed aircraft. This mid-air retrieval technique was used to recover the film canisters from the Corona spy satellites.
Types
Uncrewed
Human
The first human spaceflight was Vostok 1 on April 12, 1961, on which cosmonaut Yuri Gagarin of the USSR made one orbit around the Earth. In official Soviet documents, there is no mention of the fact that Gagarin parachuted the final seven miles. As of 2020, the only spacecraft regularly used for human spaceflight are Soyuz, Shenzhou, and Crew Dragon. The U.S. Space Shuttle fleet operated from April 1981 until July 2011. SpaceShipOne has conducted three human suborbital space flights.
Sub-orbital
On a sub-orbital spaceflight the spacecraft reaches space and then returns to the atmosphere after following a (primarily) ballistic trajectory. This is usually because of insufficient specific orbital energy, in which case a suborbital flight will last only a few minutes, but it is also possible for an object with enough energy for an orbit to have a trajectory that intersects the Earth's atmosphere, sometimes after many hours. Pioneer 1 was NASA's first space probe intended to reach the Moon. A partial failure caused it to instead follow a suborbital trajectory to an altitude of before reentering the Earth's atmosphere 43 hours after launch.
The most generally recognized boundary of space is the Kármán line above sea level. (NASA alternatively defines an astronaut as someone who has flown more than above sea level.) It is not generally recognized by the public that the increase in potential energy required to pass the Kármán line is only about 3% of the orbital energy (potential plus kinetic energy) required by the lowest possible Earth orbit (a circular orbit just above the Kármán line.) In other words, it is far easier to reach space than to stay there. On May 17, 2004, Civilian Space eXploration Team launched the GoFast rocket on a suborbital flight, the first amateur spaceflight. On June 21, 2004, SpaceShipOne was used for the first privately funded human spaceflight.
Point-to-point
Point-to-point, or Earth to Earth transportation, is a category of sub-orbital spaceflight in which a spacecraft provides rapid transport between two terrestrial locations. A conventional airline route between London and Sydney, a flight that normally lasts over twenty hours, could be traversed in less than one hour. While no company offers this type of transportation today, SpaceX has revealed plans to do so as early as the 2020s using Starship. Suborbital spaceflight over an intercontinental distance requires a vehicle velocity that is only a little lower than the velocity required to reach low Earth orbit. If rockets are used, the size of the rocket relative to the payload is similar to an Intercontinental Ballistic Missile (ICBM). Any intercontinental spaceflight has to surmount problems of heating during atmospheric re-entry that are nearly as large as those faced by orbital spaceflight.
Orbital
A minimal orbital spaceflight requires much higher velocities than a minimal sub-orbital flight, and so it is technologically much more challenging to achieve. To achieve orbital spaceflight, the tangential velocity around the Earth is as important as altitude. In order to perform a stable and lasting flight in space, the spacecraft must reach the minimal orbital speed required for a closed orbit.
Interplanetary
Interplanetary spaceflight is flight between planets within a single planetary system. In practice, the use of the term is confined to travel between the planets of our Solar System. Plans for future crewed interplanetary spaceflight missions often include final vehicle assembly in Earth orbit, such as NASA's Constellation program and Russia's Kliper/Parom tandem.
Interstellar
New Horizons is the fifth spacecraft put on an escape trajectory leaving the Solar System. Voyager 1, Voyager 2, Pioneer 10, Pioneer 11 are the earlier ones. The one farthest from the Sun is Voyager 1, which is more than 100 AU distant and is moving at 3.6 AU per year. In comparison, Proxima Centauri, the closest star other than the Sun, is 267,000 AU distant. It will take Voyager 1 over 74,000 years to reach this distance. Vehicle designs using other techniques, such as nuclear pulse propulsion are likely to be able to reach the nearest star significantly faster. Another possibility that could allow for human interstellar spaceflight is to make use of time dilation, as this would make it possible for passengers in a fast-moving vehicle to travel further into the future while aging very little, in that their great speed slows down the rate of passage of on-board time. However, attaining such high speeds would still require the use of some new, advanced method of propulsion. Dynamic soaring as a way to travel across interstellar space has been proposed as well.
Intergalactic
Intergalactic travel involves spaceflight between galaxies, and is considered much more technologically demanding than even interstellar travel and, by current engineering terms, is considered science fiction. However, theoretically speaking, there is nothing to conclusively indicate that intergalactic travel is impossible. To date several academics have studied intergalactic travel in a serious manner.
Spacecraft
Spacecraft are vehicles designed to operate in space.
The first 'true spacecraft' is sometimes said to be Apollo Lunar Module, since this was the only crewed vehicle to have been designed for, and operated only in space; and is notable for its non-aerodynamic shape.
Propulsion
Spacecraft today predominantly use rockets for propulsion, but other propulsion techniques such as ion drives are becoming more common, particularly for uncrewed vehicles, and this can significantly reduce the vehicle's mass and increase its delta-v.
Launch systems
Launch systems are used to carry a payload from Earth's surface into outer space.
Expendable
Most current spaceflight uses multi-stage expendable launch systems to reach space.
Reusable
The first reusable spacecraft, the X-15, was air-launched on a suborbital trajectory on 19 July 1963. The first partially reusable orbital spacecraft, the Space Shuttle, was launched by the USA on the 20th anniversary of Yuri Gagarin's flight, on 12 April 1981. During the Shuttle era, six orbiters were built, all of which flown in the atmosphere and five of which flown in space. The Enterprise was used only for approach and landing tests, launching from the back of a Boeing 747 and gliding to deadstick landings at Edwards AFB, California. The first Space Shuttle to fly into space was the Columbia, followed by the Challenger, Discovery, Atlantis, and Endeavour. The Endeavour was built to replace the Challenger, which was lost in January 1986. The Columbia broke up during reentry in February 2003.
The first automatic partially reusable spacecraft was the Buran (Snowstorm), launched by the USSR on 15 November 1988, although it made only one flight. This spaceplane was designed for a crew and strongly resembled the US Space Shuttle, although its drop-off boosters used liquid propellants and its main engines were located at the base of what would be the external tank in the American Shuttle. Lack of funding, complicated by the dissolution of the USSR, prevented any further flights of Buran.
The Space Shuttle was retired in 2011 due mainly to its old age. The Shuttle's human transport role is to be replaced by the SpaceX Dragon 2 and CST-100 in the 2020s. The Shuttle's heavy cargo transport role is now done by commercial launch vehicles.
Scaled Composites SpaceShipOne was a reusable suborbital spaceplane that carried pilots Mike Melvill and Brian Binnie on consecutive flights in 2004 to win the Ansari X Prize. The Spaceship Company has built its successor SpaceShipTwo. A fleet of SpaceShipTwos operated by Virgin Galactic planned to begin reusable private spaceflight carrying paying passengers (space tourists) in 2008, but this was delayed due to an accident in the propulsion development.
SpaceX achieved the first vertical soft landing of a reusable orbital rocket stage on December 21, 2015, after delivering 11 Orbcomm OG-2 commercial satellites into low Earth orbit.
The first Falcon 9 second flight occurred on 30 March 2017. SpaceX now routinely recovers and reuses their first stages, with the intent of reusing fairings as well. SpaceX is now developing a fully reusable super heavy lift rocket known as Starship, hoped to drastically reduce the price of space exploration due to being fully reusable.
Challenges
Safety
All launch vehicles contain a huge amount of energy that is needed for some part of it to reach orbit. There is therefore some risk that this energy can be released prematurely and suddenly, with significant effects. When a Delta II rocket exploded 13 seconds after launch on January 17, 1997, there were reports of store windows away being broken by the blast.
Space is a fairly predictable environment, but there are still risks of accidental depressurization and the potential failure of equipment, some of which may be very newly developed.
In April 2004 the International Association for the Advancement of Space Safety was established in the Netherlands to further international cooperation and scientific advancement in space systems safety.
Weightlessness
In a microgravity environment such as that provided by a spacecraft in orbit around the Earth, humans experience a sense of "weightlessness." Short-term exposure to microgravity causes space adaptation syndrome, a self-limiting nausea caused by derangement of the vestibular system. Long-term exposure causes multiple health issues. The most significant is bone loss, some of which is permanent, but microgravity also leads to significant deconditioning of muscular and cardiovascular tissues.
Radiation
Once above the atmosphere, radiation due to the Van Allen belts, solar radiation and cosmic radiation issues occur and increase. Further away from the Earth, solar flares can give a fatal radiation dose in minutes, and the health threat from cosmic radiation significantly increases the chances of cancer over a decade exposure or more.
Life support
In human spaceflight, the life support system is a group of devices that allow a human being to survive in outer space. NASA often uses the phrase Environmental Control and Life Support System or the acronym ECLSS when describing these systems for its human spaceflight missions. The life support system may supply: air, water and food. It must also maintain the correct body temperature, an acceptable pressure on the body and deal with the body's waste products. Shielding against harmful external influences such as radiation and micro-meteorites may also be necessary. Components of the life support system are life-critical, and are designed and constructed using safety engineering techniques.
Space weather
Space weather is the concept of changing environmental conditions in outer space. It is distinct from the concept of weather within a planetary atmosphere, and deals with phenomena involving ambient plasma, magnetic fields, radiation and other matter in space (generally close to Earth but also in interplanetary, and occasionally interstellar medium). "Space weather describes the conditions in space that affect Earth and its technological systems. Our space weather is a consequence of the behavior of the Sun, the nature of Earth's magnetic field, and our location in the Solar System."
Space weather exerts a profound influence in several areas related to space exploration and development. Changing geomagnetic conditions can induce changes in atmospheric density causing the rapid degradation of spacecraft altitude in Low Earth orbit. Geomagnetic storms due to increased solar activity can potentially blind sensors onboard spacecraft, or interfere with on-board electronics. An understanding of space environmental conditions is also important in designing shielding and life support systems for crewed spacecraft.
Environmental considerations
Exhaust pollution of rockets depends on the produced exhausts by the propellants reactions and the location of exhaustion. They mostly exhaust greenhouse gases and sometimes toxic components. Particularly at higher levels of the atmosphere the potency of exhausted gases as greenhouse gases increases considerably. Many solid rockets have chlorine in the form of perchlorate or other chemicals, and this can cause temporary local holes in the ozone layer. Re-entering spacecraft generate nitrates which also can temporarily impact the ozone layer. Most rockets are made of metals that can have an environmental impact during their construction.
While spaceflight altogether pollutes at a fraction of other human activities, it still does pollute heavily if calculated per passenger.
In addition to the atmospheric effects there are effects on the near-Earth space environment. There is the possibility that orbit could become inaccessible for generations due to exponentially increasing space debris caused by spalling of satellites and vehicles (Kessler syndrome). Many launched vehicles today are therefore designed to be re-entered after use.
Regulation
A wide range of issues such as space traffic management or liability have been issues of spaceflight regulation.
Participation and representation of all humanity in spaceflight is an issue of international space law ever since the first phase of space exploration. Even though some rights of non-spacefaring countries have been secured, sharing of space for all humanity is still criticized as imperialist and lacking, understanding spaceflight as a resource.
Access
Inclusion has been a national and international issue, resulting in 1967 in the Outer Space Treaty and its claim of outer space as the "province of all mankind". Furthermore social inclusion in human spaceflight has been demanded, with women to fly to space being limited, and minorities, like people with disability, only having been selected in European Space Agency's 2022 astronaut group.
The dominating issue about access in most recent years has been the issue of space debris and space sustainability, since established spacefaring countries endanger access to outer space with their orbital space polluting activity.
Applications
Current and proposed applications for spaceflight include:
Earth observation satellites such as spy satellites, weather satellites
Space exploration
Communication satellites
Satellite television
Satellite navigation
Space telescopes
Space tourism
Protecting Earth from potentially hazardous objects
Space colonization
Most early spaceflight development was paid for by governments. However, today major launch markets such as communication satellites and satellite television are purely commercial, though many of the launchers were originally funded by governments.
Private spaceflight is a rapidly developing area: space flight that is not only paid for by corporations or even private individuals, but often provided by private spaceflight companies. These companies often assert that much of the previous high cost of access to space was caused by governmental inefficiencies they can avoid. This assertion can be supported by much lower published launch costs for private space launch vehicles such as Falcon 9 developed with private financing. Lower launch costs and excellent safety will be required for the applications such as space tourism and especially space colonization to become feasible for expansion.
Spacefaring
To be spacefaring is to be capable of and active in the operation of spacecraft. It involves a knowledge of a variety of topics and development of specialised skills including: aeronautics; astronautics; programs to train astronauts; space weather and forecasting; spacecraft operations; operation of various equipment; spacecraft design and construction; atmospheric takeoff and reentry; orbital mechanics (a.k.a. astrodynamics); communications; engines and rockets; execution of evolutions such as towing, microgravity construction, and space docking; cargo handling equipment, dangerous cargos and cargo storage; spacewalking; dealing with emergencies; survival at space and first aid; fire fighting; life support. The degree of knowledge needed within these areas is dependent upon the nature of the work and the type of vessel employed. "Spacefaring" is analogous to seafaring.
There has never been a crewed mission outside the Earth–Moon system. However, the United States, Russia, China, European Space Agency (ESA) countries, and a few corporations and enterprises have plans in various stages to travel to Mars (see Human mission to Mars).
Spacefaring entities can be sovereign states, supranational entities, and private corporations. Spacefaring nations are those capable of independently building and launching craft into space. A growing number of private entities have become or are becoming spacefaring.
Global coordination
The United Nations Office for Outer Space Affairs (UNOOSA) has been the main multilateral body servicing international contact and exchange on space activity among spacefaring and non-spacefaring states.
Crewed spacefaring nations
Currently Russia, the United States and China are the only crewed spacefaring nations.
Spacefaring nations listed by date of first crewed launch:
Soviet Union (Russia) (1961)
United States (1961)
China (2003)
Uncrewed spacefaring nations
The following nations or organizations have developed their own launch vehicles to launch uncrewed spacecraft into orbit either from their own territory or with foreign assistance (date of first launch in parentheses):
Soviet Union (1957)
United States (1958)
France (1965)
Italy (1967)★
Australia (1967)★
Japan (1970)
China (1970)
United Kingdom (1971)
European Space Agency (1979)
India (1980)
Israel (1988)
Ukraine (1991)*
Russia (1992)*
Iran (2009)
North Korea (2012)
South Korea (2013)★
New Zealand (2018)★
*Previously a major region in the Soviet Union
★Launch vehicle fully or partially developed by another country
Also several countries, such as Canada, Italy, and Australia, had semi-independent spacefaring capability, launching locally-built satellites on foreign launchers. Canada had designed and built satellites (Alouette 1 and 2) in 1962 and 1965 which were orbited using U.S. launch vehicles. Italy has designed and built several satellites, as well as pressurized modules for the International Space Station. Early Italian satellites were launched using vehicles provided by NASA, first from Wallops Flight Facility in 1964 and then from a spaceport in Kenya (San Marco Platform) between 1967 and 1988; Italy has led the development of the Vega rocket programme within the European Space Agency since 1998. The United Kingdom abandoned its independent space launch program in 1972 in favour of co-operating with the European Launcher Development Organisation (ELDO) on launch technologies until 1974. Australia abandoned its launcher program shortly after the successful launch of WRESAT, and became the only non-European member of ELDO.
Suborbital
Considering merely launching an object beyond the Kármán line to be the minimum requirement of spacefaring, Germany, with the V-2 rocket, became the first spacefaring nation in 1944. The following nations have only achieved suborbital spaceflight capability by launching indigenous rockets or missiles or both into suborbital space:
Nazi Germany (June 20, 1944)
East Germany (April 12, 1957)
Canada (September 5, 1959)
Lebanon (November 21, 1962)
Switzerland (October 27, 1967)
Argentina (April 16, 1969)
Brazil (September 21, 1976)
Spain (February 18, 1981)
West Germany (March 1, 1981)
Iraq (June 1984)
South Africa (June 1, 1989)
Sweden (May 8, 1991)
Yemen (May 12, 1994)
Pakistan (April 6, 1998)
Taiwan (December 15, 1998)
Syria (September 1, 2000)
Indonesia (September 29, 2004)
Democratic Republic of the Congo (2007)
New Zealand (November 30, 2009)
Norway (September 27, 2018)
Netherlands (September 19, 2020)
Turkey (October 29, 2020)
| Technology | Space | null |
77229 | https://en.wikipedia.org/wiki/Microsoft%20PowerPoint | Microsoft PowerPoint | Microsoft PowerPoint is a presentation program, created by Robert Gaskins, Tom Rudkin, and Dennis Austin at a software company named Forethought, Inc. It was released on April 20, 1987, initially for Macintosh computers only. Microsoft acquired PowerPoint for about $14 million three months after it appeared. This was Microsoft's first significant acquisition, and Microsoft set up a new business unit for PowerPoint in Silicon Valley where Forethought had been located.
PowerPoint became a component of the Microsoft Office suite, first offered in 1989 for Macintosh and in 1990 for Windows, which bundled several Microsoft apps. Beginning with PowerPoint 4.0 (1994), PowerPoint was integrated into Microsoft Office development, and adopted shared common components and a converged user interface.
PowerPoint's market share was very small at first, prior to introducing a version for Microsoft Windows, but grew rapidly with the growth of Windows and of Office. Since the late 1990s, PowerPoint's worldwide market share of presentation software has been estimated at 95 percent.
PowerPoint was originally designed to provide visuals for group presentations within business organizations, but has come to be widely used in other communication situations in business and beyond. The wider use led to the development of the PowerPoint presentation as a new form of communication, with strong reactions including advice that it should be used less, differently, or better.
The first PowerPoint version (Macintosh 1987) was used to produce overhead transparencies, the second (Macintosh 1988, Windows 1990) could also produce color 35 mm slides. The third version (Windows and Macintosh 1992) introduced video output of virtual slideshows to digital projectors, which would over time replace physical transparencies and slides. A dozen major versions since then have added additional features and modes of operation and have made PowerPoint available beyond Apple Macintosh and Microsoft Windows, adding versions for iOS, Android, and web access.
History
Creation at Forethought (1984–1987)
PowerPoint was created by Robert Gaskins and Dennis Austin at a software startup in Silicon Valley named Forethought, Inc. Forethought had been founded in 1983 to create an integrated environment and applications for future personal computers that would provide a graphical user interface, but it had run into difficulties requiring a "restart" and new plan.
On July 5, 1984, Forethought hired Robert Gaskins as its vice president of product development to create a new application that would be especially suited to the new graphical personal computers, such as the Apple Macintosh and later Microsoft Windows. Gaskins produced his initial description of PowerPoint about a month later (August 14, 1984) in the form of a 2-page document titled "Presentation Graphics for Overhead Projection." By October 1984, Gaskins had selected Dennis Austin to be the developer for PowerPoint. Gaskins and Austin worked together on the definition and design of the new product for nearly a year, and produced the first specification document dated August 21, 1985. This first design document showed a product as it would look in Microsoft Windows 1.0, which at that time had not been released.
Development from that spec was begun by Austin in November 1985, for Macintosh first. About six months later, on May 1, 1986, Gaskins and Austin chose a second developer to join the project, Thomas Rudkin. Gaskins prepared two final product specification marketing documents in June 1986; these described a product for both Macintosh and Windows. At about the same time, Austin, Rudkin, and Gaskins produced a second and final major design specification document, this time showing a Macintosh look.
Throughout this development period, the product was called "Presenter". Then, just before release, there was a last-minute check with Forethought's lawyers to register the name as a trademark, and "Presenter" was unexpectedly rejected because it had already been used by someone else. Gaskins says that he thought of "PowerPoint", based on the product's goal of "empowering" individual presenters, and sent that name to the lawyers for clearance, while all the documentation was hastily revised.
Funding to complete development of PowerPoint was assured in mid-January 1987, when a new Apple Computer venture capital fund, called Apple's Strategic Investment Group, selected PowerPoint to be its first investment. A month later, on February 22, 1987, Forethought announced PowerPoint at the Personal Computer Forum in Phoenix; John Sculley, the CEO of Apple, appeared at the announcement and said "We see desktop presentation as potentially a bigger market for Apple than desktop publishing."
PowerPoint 1.0 for Macintosh shipped from manufacturing on April 20, 1987, and the first production run of 10,000 units was sold out.
Acquisition by Microsoft (1987–1992)
By early 1987, Microsoft was starting to plan a new application to create presentations, an activity led by Jeff Raikes, who was head of marketing for the Applications Division. Microsoft assigned an internal group to write a specification and plan for a new presentation product. They contemplated an acquisition to speed up development, and in early 1987 Microsoft sent a letter of intent to acquire Dave Winer's product called MORE, an outlining program that could print its outlines as bullet charts. During this preparatory activity Raikes discovered that a program specifically to make overhead presentations was already being developed by Forethought, Inc., and that it was nearly completed. Raikes and others visited Forethought on February 6, 1987, for a confidential demonstration.
Raikes later recounted his reaction to seeing PowerPoint and his report about it to Bill Gates, who was initially skeptical:
When PowerPoint was released by Forethought, its initial press was favorable; the Wall Street Journal reported on early reactions: I see about one product a year I get this excited about,' says Amy hora, a consultant in Bala Cynwyd, Pa. 'People will buy a Macintosh just to get access to this product.
On April 28, 1987, a week after shipment, a group of Microsoft's senior executives spent another day at Forethought to hear about initial PowerPoint sales on Macintosh and plans for Windows. The following day, Microsoft sent a letter to Dave Winer withdrawing its earlier letter of intent to acquire his company, and in mid-May 1987 Microsoft sent a letter of intent to acquire Forethought. As requested in that letter of intent, Robert Gaskins from Forethought went to Redmond for a one-on-one meeting with Bill Gates in early June 1987, and by the end of July an agreement was concluded for an acquisition. The New York Times reported:
Microsoft's president Jon Shirley offered his company's motivation for the acquisition: We made this deal primarily because of our belief in desktop presentations as a product category. ... Forethought was first to market with a product in this category. Microsoft had 50% market share in Macintosh applications, and led in three categories; Raikes said that after the acquisition it would lead in five categories. (Forethought distributed the database Filemaker, which Microsoft wanted to continue marketing.) The company intended for Forethought to be its Silicon Valley base to develop and market future graphics software, so set up within its Applications Division, an independent "Graphics Business Unit" for PowerPoint, the first Microsoft application group distant from the main Redmond location. The company hoped to hire employees uninterested in living in Washington state. All the PowerPoint people from Forethought joined Microsoft, and the new location was headed by Robert Gaskins, with Dennis Austin and Thomas Rudkin leading development. PowerPoint 1.0 for Macintosh was modified to indicate the new Microsoft ownership and continued to be sold. A year after the acquisition, Gaskins reported that all seven Forethought PowerPoint employees had stayed with Microsoft, and the Graphics Business Unit had hired 12 employees, many of whom did not want to move to Redmond. The GBU had moved to a new location on Sand Hill Road in Menlo Park, California; it was much larger than needed for 19 people, but Gaskins wrote that he and Microsoft wanted future capacity as the company grew in Silicon Valley.
A new PowerPoint 2.0 for Macintosh, adding color 35 mm slides, shipped in May 1988, and again received good reviews. The same PowerPoint 2.0 product re-developed for Windows was shipped two years later, in mid-1990, at the same time as Windows 3.0. Much of the color technology was the result of a joint development partnership with Genigraphics, the dominant presentation services company.
PowerPoint 3.0, which was shipped in 1992 for both Windows and Mac, added live video for projectors and monitors, with the result that PowerPoint was thereafter used for delivering presentations as well as for preparing them. This was at first an alternative to overhead transparencies and 35 mm slides, but over time would come to replace them.
Part of Microsoft Office (since 1993)
PowerPoint had been included in Microsoft Office from the beginning. PowerPoint 2.0 for Macintosh was part of the first Office bundle for Macintosh which was offered in mid-1989. When PowerPoint 2.0 for Windows appeared, a year later, it was part of a similar Office bundle for Windows, which was offered in late 1990. Both of these were bundling promotions, in which the independent applications were packaged together and offered for a lower total price.
PowerPoint 3.0 (1992) was again separately specified and developed, and was advertised and sold separately from Office. It was, as before, included in Microsoft Office 3.0, both for Windows and the corresponding version for Macintosh.
A plan to integrate the applications themselves more tightly had been indicated as early as February 1991, toward the end of PowerPoint 3.0 development, in an internal memo by Bill Gates:
The move from bundling separate products to integrated development began with PowerPoint 4.0, developed in 1993–1994 under new management from Redmond. The PowerPoint group in Silicon Valley was reorganized from the independent "Graphics Business Unit" (GBU) to become the "Graphics Product Unit" (GPU) for Office, and PowerPoint 4.0 changed to adopt a converged user interface and other components shared with the other apps in Office.
When it was released, the computer press reported on the change approvingly: "PowerPoint 4.0 has been re-engineered from the ground up to resemble and work with the latest applications in Office: Word 6.0, Excel 5.0, and Access 2.0. The integration is so good, you'll have to look twice to make sure you're running PowerPoint and not Word or Excel." Office integration was further underscored in the following version, PowerPoint 95, which was given the version number PowerPoint 7.0 (skipping 5.0 and 6.0) so that all the components of Office would share the same major version number.
Although PowerPoint by this point had become part of the integrated Microsoft Office product, its development remained in Silicon Valley. Succeeding versions of PowerPoint introduced important changes, particularly version 12.0 (2007) which had a very different shared Office "ribbon" user interface, and a new shared Office XML-based file format. This marked the 20th anniversary of PowerPoint, and Microsoft held an event to commemorate that anniversary at its Silicon Valley Campus for the PowerPoint team there. Special guests were Robert Gaskins, Dennis Austin, and Thomas Rudkin, and the featured speaker was Jeff Raikes, all from PowerPoint 1.0 days, 20 years before.
Since then major development of PowerPoint as part of Office has continued. New development techniques (shared across Office) for PowerPoint 2016 have made it possible to ship versions of PowerPoint 2016 for Windows, Mac, iOS, Android, and web access nearly simultaneously, and to release new features on an almost monthly schedule. PowerPoint development is still carried out in Silicon Valley .
In 2010, Jeff Raikes, who had most recently been President of the Business Division of Microsoft (including responsibility for Office), observed: "of course, today we know that PowerPoint is oftentimes the number two—or in some cases even the number one—most-used tool" among the applications in Office.
Sales and market share
PowerPoint's initial sales were about 40,000 copies sold in 1987 (nine months), about 85,000 copies in 1988, and about 100,000 copies in 1989, all for Macintosh. PowerPoint's market share in its first three years was a tiny part of the total presentation market, which was very heavily dominated by MS-DOS applications on PCs. The market leaders on MS-DOS in 1988–1989 were Harvard Graphics (introduced by Software Publishing in 1986) in first place, and Lotus Freelance Plus (also introduced in 1986) as a strong second. They were competing with more than a dozen other MS-DOS presentation products, and Microsoft did not develop a PowerPoint version for MS-DOS. After three years, PowerPoint sales were disappointing. Jeff Raikes, who had bought PowerPoint for Microsoft, later recalled: "By 1990, it looked like it wasn't a very smart idea [for Microsoft to have acquired PowerPoint], because not very many people were using PowerPoint."
This began to change when the first version for Windows, PowerPoint 2.0, brought sales up to about 200,000 copies in 1990 and to about 375,000 copies in 1991, with Windows units outselling Macintosh. PowerPoint sold about 1 million copies in 1992, of which about 80 percent were for Windows and about 20 percent for Macintosh, and in 1992 PowerPoint's market share of worldwide presentation graphics software sales was reported as 63 percent. By the last six months of 1992, PowerPoint revenue was running at a rate of over $100 million annually ($ in present-day terms).
Sales of PowerPoint 3.0 doubled to about 2 million copies in 1993, of which about 90 percent were for Windows and about 10 percent for Macintosh, and in 1993 PowerPoint's market share of worldwide presentation graphics software sales was reported as 78 percent. In both years, about half of total revenue came from sales outside the U.S.
By 1997 PowerPoint sales had doubled again, to more than 4 million copies annually, representing 85 percent of the world market. Also in 1997, an internal publication from the PowerPoint group said that by then over 20 million copies of PowerPoint were in use, and that total revenues from PowerPoint over its first ten years (1987 to 1996) had already exceeded $1 billion.
Since the late 1990s, PowerPoint's market share of total world presentation software has been estimated at 95 percent by both industry and academic sources.
Operation
The earliest version of PowerPoint (1987 for Macintosh) could be used to print black and white pages to be photocopied onto sheets of transparent film for projection from overhead projectors, and to print speaker's notes and audience handouts; the next version (1988 for Macintosh, 1990 for Windows) was extended to also produce color 35mm slides by communicating a file over a modem to a Genigraphics imaging center with slides returned by overnight delivery for projection from slide projectors. PowerPoint was used for planning and preparing a presentation, but not for delivering it (apart from previewing it on a computer screen, or distributing printed paper copies). The operation of PowerPoint changed substantially in its third version (1992 for Windows and Macintosh), when PowerPoint was extended to also deliver a presentation by producing direct video output to digital projectors or large monitors. In 1992 video projection of presentations was rare and expensive, and practically unknown from a laptop computer. Robert Gaskins, one of the creators of PowerPoint, says he publicly demonstrated that use for the first time at a large Microsoft meeting held in Paris on February 25, 1992, by using an unreleased development build of PowerPoint 3.0 running on an early pre-production sample of a powerful new color laptop and feeding a professional auditorium video projector.
By about 2003, ten years later, digital projection had become the dominant mode of use, replacing transparencies and 35mm slides and their projectors. As a result, the meaning of "PowerPoint presentation" narrowed to mean specifically digital projection:
In contemporary operation, PowerPoint is used to create a file (called a "presentation" or "deck") containing a sequence of pages (called "slides" in the app) which usually have a consistent style (from template masters), and which may contain information imported from other apps or created in PowerPoint, including text, bullet lists, tables, charts, drawn shapes, images, audio clips, video clips, animations of elements, and animated transitions between slides, plus attached notes for each slide.
After such a file is created, typical operation is to present it as a slide show using a portable computer, where the presentation file is stored on the computer or available from a network, and the computer's screen shows a "presenter view" with current slide, next slide, speaker's notes for the current slide, and other information. Video is sent from the computer to one or more external digital projectors or monitors, showing only the current slide to the audience, with sequencing controlled by the speaker at the computer. A smartphone remote control built in to PowerPoint for iOS (optionally controlled from Apple Watch) and for Android allows the presenter to control the show from elsewhere in the room.
In addition to a computer slide show projected to a live audience by a speaker, PowerPoint can be used to deliver a presentation in a number of other ways:
Displayed on the screen of the presentation computer or tablet (for a very small group)
Printed for distribution as paper documents (in several formats)
Distributed as files for private viewing, even on computers without PowerPoint
Packaged for distribution on CD or a network, including linked and embedded data
Transmitted as a live broadcast presentation over the web
Embedded in a web page or blog
Shared on social networks such as Facebook or Twitter
Set up as a self-running unattended display
Recorded as video/audio (H.264/AAC), to be distributed as for any other video
Some of these ways of using PowerPoint have been studied by JoAnne Yates and Wanda Orlikowski of the MIT Sloan School of Management:
They found that some of these ways of using PowerPoint could influence the content of presentations, for example when "the slides themselves have to carry more of the substance of the presentation, and thus need considerably more content than they would have if they were intended for projection by a speaker who would orally provide additional details and nuance about content and context."
Other platforms
PowerPoint for mobile
PowerPoint Mobile is included with Windows Mobile 5.0. It is a presentation program capable of reading and editing Microsoft PowerPoint presentations, although authoring abilities are limited to adding notes, editing text, and rearranging slides. It can't create new presentations. Versions of PowerPoint Mobile for Windows Phone 7 can also watch presentation broadcasts streamed from the Internet. In 2015, Microsoft released PowerPoint Mobile for Windows 10 as a universal app. In this version of PowerPoint users can create and edit new presentations, present, and share their PowerPoint documents.
PowerPoint for the web
PowerPoint for the web is a free lightweight version of Microsoft PowerPoint available as part of Office on the web, which also includes web versions of Microsoft Excel and Microsoft Word.
PowerPoint for the web does not support inserting or editing charts, equations, or audio or video stored on your PC, but they are all displayed in the presentation if they were added in using a desktop app. Some elements, like WordArt effects or more advanced animations and transitions, are not displayed at all, although they are preserved in the document. PowerPoint for the web also lacks the Outline, Master, Slide Sorter, and Presenter views present in the desktop app, as well as having limited printing options.
Cultural impact
Business uses
PowerPoint was originally targeted just for business presentations. Robert Gaskins, who was responsible for its design, has written about his intended customers: "... I did not target other existing large groups of users of presentations, such as school teachers or military officers. ... I also did not plan to target people who were not existing users of presentations ... such as clergy and school children ... . Our focus was purely on business users, in small and large companies, from one person to the largest multinationals." Business people had for a long time made presentations for sales calls and for internal company communications, and PowerPoint produced the same formats in the same style and for the same purposes.
PowerPoint use in business grew over its first five years (1987–1992) to sales of about 1 million copies annually, for worldwide market share of 63 percent. Over the following five years (1992–1997) PowerPoint sales accelerated, to a rate of about 4 million copies annually, for worldwide market share of 85 percent. The increase in business use has been attributed to "network effects", whereby additional users of PowerPoint in a company or an industry increased its salience and value to other users.
Not everyone immediately approved of the greater use of PowerPoint for presentations, even in business. CEOs who very early were reported to discourage or ban PowerPoint presentations at internal business meetings included Lou Gerstner (at IBM, in 1993), Scott McNealy (at Sun Microsystems, in 1996), and Steve Jobs (at Apple, in 1997). But even so, Rich Gold, a scholar who studied corporate presentation use at Xerox PARC, could write in 1999: "Within today's corporation, if you want to communicate an idea ... you use PowerPoint."
Uses beyond business
At the same time that PowerPoint was becoming dominant in business settings, it was also being adopted for uses beyond business: "Personal computing ... scaled up the production of presentations. ... The result has been the rise of presentation culture. In an information society, nearly everyone presents."
In 1998, at about the same time that Gold was pronouncing PowerPoint's ubiquity in business, the influential Bell Labs engineer Robert W. Lucky could already write about broader uses:
Over a decade or so, beginning in the mid-1990s, PowerPoint began to be used in many communication situations, well beyond its original business presentation uses, to include teaching in schools and in universities, lecturing in scientific meetings (and preparing their related poster sessions), worshipping in churches, making legal arguments in courtrooms, displaying supertitles in theaters, driving helmet-mounted displays in spacesuits for NASA astronauts, giving military briefings, issuing governmental reports, undertaking diplomatic negotiations, writing novels, giving architectural demonstrations, prototyping website designs, creating animated video games, editing images, creating art projects, and even as a substitute for writing engineering technical reports, and as an organizing tool for writing general business documents.
By 2003, it seemed that PowerPoint was being used everywhere. Julia Keller reported for the Chicago Tribune:
Cultural reactions
As uses broadened, cultural awareness of PowerPoint grew and commentary about it began to appear. "With the widespread adoption of PowerPoint came complaints ... often very general statements reflecting dissatisfaction with modern media and communication practices as well as the dysfunctions of organizational culture." Indications of this awareness included increasing mentions of PowerPoint use in the Dilbert comic strips of Scott Adams, comic parodies of poor or inappropriate use such as the Gettysburg Address in PowerPoint or summaries of Shakespeare's Hamlet and Nabokov's Lolita in PowerPoint, and a vast number of publications on the general subject of PowerPoint, especially about how to use it.
Out of all the analyses of PowerPoint over a quarter of a century, at least three general themes emerged as categories of reaction to its broader use: (1) "Use it less": avoid PowerPoint in favor of alternatives, such as using more-complex graphics and written prose, or using nothing; (2) "Use it differently": make a major change to a PowerPoint style that is simpler and pictorial, turning the presentation toward a performance, more like a Steve Jobs keynote; and (3) "Use it better": retain much of the conventional PowerPoint style but learn to avoid making many kinds of mistakes that can interfere with communication.
Use it less
An early reaction was that the broader use of PowerPoint was a mistake, and should be reversed. An influential example of this came from Edward Tufte, an authority on information design, who has been a professor of political science, statistics, and computer science at Princeton and Yale, but is best known for his self-published books on data visualization, which have sold nearly 2 million copies as of 2014.
In 2003, he published a widely-read booklet titled The Cognitive Style of PowerPoint, revised in 2006. Tufte found a number of problems with the "cognitive style" of PowerPoint, many of which he attributed to the standard default style templates:
Tufte particularly advised against using PowerPoint for reporting scientific analyses, using as a dramatic example some slides made during the flight of the space shuttle Columbia after it had been damaged by an accident at liftoff, slides which poorly communicated the engineers' limited understanding of what had happened. For such technical presentations, and for most occasions apart from its initial domain of sales presentations, Tufte advised against using PowerPoint at all; in many situations, according to Tufte, it would be better to substitute high-resolution graphics or concise prose documents as handouts for the audience to study and discuss, providing a great deal more detail.
Many commentators enthusiastically joined in Tufte's vivid criticism of PowerPoint uses, and at a conference held in 2013 (a decade after Tufte's booklet appeared) one paper claimed that "Despite all the criticism about his work, Tufte can be considered as the single most influential author in the discourse on PowerPoint. ... While his approach was not rigorous from a research perspective, his articles received wide resonance with the public at large ... ." There were also others who disagreed with Tufte's assertion that the PowerPoint program reduces the quality of presenters' thoughts: Steven Pinker, professor of psychology at MIT and later Harvard, had earlier argued that "If anything, PowerPoint, if used well, would ideally reflect the way we think." Pinker later reinforced this opinion: "Any general opposition to PowerPoint is just dumb, ... It's like denouncing lectures—before there were awful PowerPoint presentations, there were awful scripted lectures, unscripted lectures, slide shows, chalk talks, and so on."
Much of the early commentary, on all sides, was "informal" and "anecdotal", because empirical research had been limited.
Use it differently
A second reaction to PowerPoint use was to say that PowerPoint can be used well, but only by substantially changing its style of use. This reaction is exemplified by Richard E. Mayer, a professor of psychology at the University of California, Santa Barbara, who has studied cognition and learning, particularly the design of educational multimedia, and who has published more than 500 publications, including over 30 books. Mayer's theme has been that "In light of the science, it is up to us to make a fundamental shift in our thinking—we can no longer expect people to struggle to try to adapt to our PowerPoint habits. Instead, we have to change our PowerPoint habits to align with the way people learn."
Tufte had argued his judgment that the information density of text on PowerPoint slides was too low, perhaps only 40 words on a slide, leading to over-simplified messages; Mayer responded that his empirical research showed exactly the opposite, that the amount of text on PowerPoint slides was usually too high, and that even fewer than 40 words on a slide resulted in "PowerPoint overload" that impeded understanding during presentations.
Mayer suggested a few major changes from traditional PowerPoint formats:
replacing brief slide titles with longer "headlines" expressing complete ideas;
showing more slides but simpler ones;
removing almost all text including nearly all bullet lists (reserving the text for the spoken narration);
using larger, higher-quality, and more important graphics and photographs;
removing all extraneous decoration, backgrounds, logos and identifications, everything but the essential message.
Mayer's ideas are claimed by Carmine Gallo to have been reflected in Steve Jobs's presentations: "Mayer outlined fundamental principles of multimedia design based on what scientists know about cognitive functioning. Steve Jobs's slides adhere to each of Mayer's principles ... ." Though not unique to Jobs, many people saw the style for the first time in Jobs's famous product introductions. Steve Jobs would have been using Apple's Keynote, which was designed for Jobs's own slide shows beginning in 2003, but Gallo says that "speaking like Jobs has little to do with the type of presentation software you use (PowerPoint, Keynote, etc.) ... all the techniques apply equally to PowerPoint and Keynote." Gallo adds that "Microsoft's PowerPoint has one big advantage over Apple's Keynote presentation software—it's everywhere ... it's safe to say that the number of Keynote presentations is minuscule in comparison with PowerPoint. Although most presentation designers who are familiar with both formats prefer to work in the more elegant Keynote system, those same designers will tell you that the majority of their client work is done in PowerPoint."
Consistent with its association with Steve Jobs's keynotes, a response to this style has been that it is particularly effective for "ballroom-style presentations" (as often given in conference center ballrooms) where a celebrated and practiced speaker addresses a large passive audience, but less appropriate for "conference room-style presentations" which are often recurring internal business meetings for in-depth discussion with motivated counterparts.
Use it better
A third reaction to PowerPoint use was to conclude that the standard style is capable of being used well, but that many small points need to be executed carefully, to avoid impeding understanding. This kind of analysis is particularly associated with Stephen Kosslyn, a cognitive neuroscientist who specializes in the psychology of learning and visual communication, and who has been head of the department of psychology at Harvard, has been Director of Stanford's Center for Advanced Study in the Behavioral Sciences, and has published some 300 papers and 14 books.
Kosslyn presented a set of psychological principles of "human perception, memory, and comprehension" that "appears to capture the major points of agreement among researchers." He reports that his experiments support the idea that it is not intuitive or obvious how to create effective PowerPoint presentations that conform to those agreed principles, and that even small differences that might not seem significant to a presenter can produce very different results in audiences' understanding. For this reason, Kosslyn says, users need specific education to be able to identify best ways to avoid "flaws and failures":
The many "flaws and failures" identified were those "likely to disrupt the comprehension or memory of the material." Among the most common examples were "Bulleted items are not presented individually, growing the list from the top to the bottom," "More than four bulleted items appear in a single list," "More than two lines are used per bulleted sentence," and "Words are not large enough (i.e., greater than 20 point) to be easily seen." Among audience reactions common problems reported were "Speakers read word-for-word from notes or from the slides themselves," "The slides contained too much material to absorb before the next slide was presented," and "The main point was obscured by lots of irrelevant detail."
Kosslyn observes that these findings could help to explain why the many studies of the instructional effectiveness of PowerPoint have been inconclusive and conflicting, if there were differences in the quality of the presentations tested in different studies that went unobserved because "many may feel that 'good design' is intuitively clear."
In 2007 Kosslyn wrote a book about PowerPoint, in which he suggested a very large number of fairly modest changes to PowerPoint styles and gave advice on recommended ways of using PowerPoint. In a later second book about PowerPoint he suggested nearly 150 clarifying style changes (in fewer than 150 pages). Kosslyn summarizes:
In 2017, an online poll of social media users in the UK was reported to show that PowerPoint "remains as popular with young tech-savvy users as it is with the Baby Boomers," with about four out of five saying that "PowerPoint was a great tool for making presentations," in part because "PowerPoint, with its capacity to be highly visual, bridges the wordy world of yesterday with the visual future of tomorrow."
Also in 2017, the Managerial Communication Group of MIT Sloan School of Management polled their incoming MBA students, finding that "results underscore just how differently this generation communicates as compared with older workers." Fewer than half of respondents reported doing any meaningful, longer-form writing at work, and even that minority mostly did so very infrequently, but "85 percent of students named producing presentations as a meaningful part of their job responsibilities. Two-thirds report that they present on a daily or weekly basis—so it's no surprise that in-person presentations is the top skill they hope to improve." One of the researchers concluded: "We're not likely to see future workplaces with long-form writing. The trend is toward presentations and slides, and we don't see any sign of that slowing down."
U.S. military excess
Use of PowerPoint by the U.S. military services began slowly, because they were invested in mainframe computers, MS-DOS PCs and specialized military-specification graphic output devices, all of which PowerPoint did not support. But because of the strong military tradition of presenting briefings, as soon as they acquired the computers needed to run it, PowerPoint became part of the U.S. military.
By 2000, ten years after PowerPoint for Windows appeared, it was already identified as an important feature of U.S. armed forces culture, in a front-page story in the Wall Street Journal:
U.S. military use of PowerPoint may have influenced its use by armed forces of other countries: "Foreign armed services also are beginning to get in on the act. 'You can't speak with the U.S. military without knowing PowerPoint,' says Margaret Hayes, an instructor at National Defense University in Washington D.C., who teaches Latin American military officers how to use the software."
After another 10 years, in 2010 (and again on its front page) the New York Times reported that PowerPoint use in the military was then "a military tool that has spun out of control":
The New York Times account went on to say that as a result some U.S. generals had banned the use of PowerPoint in their operations:
Several incidents, about the same time, gave wide currency to discussions by serving military officers describing excessive PowerPoint use and the organizational culture that encouraged it. In response to the New York Times story, Peter Norvig and Stephen M. Kosslyn sent a joint letter to the editor stressing the institutional culture of the military: "... many military personnel bemoan the overuse and misuse of PowerPoint. ... The problem is not in the tool itself, but in the way that people use it—which is partly a result of how institutions promote misuse."
The two generals who had been mentioned in 2010 as opposing the institutional culture of excessive PowerPoint use were both in the news again in 2017, when James N. Mattis became U.S. Secretary of Defense, and H. R. McMaster was appointed as U.S. National Security Advisor.
Artistic medium
Musician David Byrne has been using PowerPoint as a medium for art for years, producing a book and DVD and showing at galleries his PowerPoint-based artwork. Byrne has written: "I have been working with PowerPoint, the ubiquitous presentation software, as an art medium for a number of years. It started off as a joke (this software is a symbol of corporate salesmanship, or lack thereof) but then the work took on a life of its own as I realized I could create pieces that were moving, despite the limitations of the 'medium.
In 2005 Byrne toured with a theater piece styled as a PowerPoint presentation. When he presented it in Berkeley, on March 8, 2005, the University of California news service reported: "Byrne also defended [PowerPoint's] appeal as more than just a business tool—as a medium for art and theater. His talk was titled 'I ♥ PowerPoint'. Berkeley alumnus Bob Gaskins and Dennis Austin were in the audience. Eventually, Byrne said, PowerPoint could be the foundation for 'presentational theater,' with roots in Brechtian drama and Asian puppet theater." After that performance, Byrne described it in his own online journal: "Did the PowerPoint talk in Berkeley for an audience of IT legends and academics. I was terrified. The guys that originally turned PowerPoint into a program were there, what were THEY gonna think? ... [Gaskins] did tell me afterwards that he liked the PowerPoint as theater idea, which was a relief."
The expressions "PowerPoint Art" or "pptArt" are used to define a contemporary Italian artistic movement which believes that the corporate world can be a unique and exceptional source of inspiration for the artist. They say: "The pptArt name refers to PowerPoint, the symbolic and abstract language developed by the corporate world which has become a universal and highly symbolic communication system beyond cultures and borders."
The wide use of PowerPoint had, by 2010, given rise to " ... a subculture of PowerPoint enthusiasts [that] is teaching the old application new tricks, and may even be turning a dry presentation format into a full-fledged artistic medium," by using PowerPoint animation to create "games, artworks, anime, and movies."
PowerPoint Viewer
PowerPoint Viewer is the name for a series of small free application programs to be used on computers without PowerPoint installed, to view, project, or print (but not create or edit) presentations.
The first version was introduced with PowerPoint 3.0 in 1992, to enable electronic presentations to be projected using conference-room computers and to be freely distributed; on Windows, it took advantage of the new feature of embedding TrueType fonts within PowerPoint presentation files to make such distribution easier. The same kind of viewer app was shipped with PowerPoint 3.0 for Macintosh, also in 1992.
Beginning with PowerPoint 2003, a feature called "Package for CD" automatically managed all linked video and audio files plus needed fonts when exporting a presentation to a disk or flash drive or network location, and also included a copy of a revised PowerPoint Viewer application so that the result could be presented on other PCs without installing anything.
The latest version that runs on Windows "was created in conjunction with PowerPoint 2010, but it can also be used to view newer presentations created in PowerPoint 2013 and PowerPoint 2016. ... All transitions, videos and effects appear and behave the same when viewed using PowerPoint Viewer as they do when viewed in PowerPoint 2010." It supports presentations created using PowerPoint 97 and later. The latest version that runs on Macintosh is PowerPoint 98 Viewer for the Classic Mac OS and Classic Environment, for Macs supporting System 7.5 to Mac OS X Tiger (10.4). It can open presentations only from PowerPoint 3.0, 4.0, and 8.0 (PowerPoint 98), although presentations created on Mac can be opened in PowerPoint Viewer on Windows.
, the last versions of PowerPoint Viewer for all platforms have been retired by Microsoft; they are no longer available for download and no longer receive security updates. The final PowerPoint Viewer for Windows (2010) and the final PowerPoint Viewer for Classic Mac OS (1998) are available only from archives. The recommended replacements for PowerPoint Viewer: "On Windows 10 PCs, download the free ... PowerPoint Mobile application from the Windows Store," and "On Windows 7 or Windows 8/8.1 PCs, upload the file to OneDrive and view it for free using ... PowerPoint Online."
Versions
PowerPoint 1.0
For Macintosh: April 1987
Innovations included: multiple slides in a single file, organizing slides with a slide sorter view and a title view (precursor of outline view), speakers' notes pages attached to each slide, printing of audience handouts with multiple slides per page, text with outlining styles and full word-processor formatting, graphic shapes with attached text for drawing diagrams and tables. It also shipped with a hardbound book as its manual.
"It produced overhead transparencies on a black-and-white Macintosh for laser printing. Presenters could now directly control their own overheads and would no longer have to work through the person with the typewriter. PowerPoint handled the task of making the overheads all look alike; one change reformats them all. Typographic fonts were better than an Orator typeball, and charts and diagrams could be imported from MacDraw, MacPaint, and Excel, thanks to the new Mac clipboard."
System requirements: (Mac) Original Macintosh or better, System 1.0 or higher, 512K RAM.
PowerPoint 2.0
For Macintosh: May 1988; for Windows: May 1990
Part of Microsoft Office for Mac and Microsoft Office for Windows. Innovations included: color, more word processing features, find and replace, spell checking, color schemes for presentations, guide to color selection, ability to change color scheme retrospectively, shaded coloring for fills.
"It added color 35 mm slides, transmitting the resulting file over a modem to Genigraphics for imaging on Genigraphics' film recorders and photo processing in Genigraphics' labs overnight. Genigraphics was the leading professional service bureau, having developed its own Digital Equipment Corp. PDP-11-based computer systems for its artists. After a short time, though, Genigraphics itself switched to PowerPoint."
System requirements: (Mac) Original Macintosh or better, System 4.1 or higher, 1 MB RAM. (Windows) 286 PC or higher, Windows 3.0, 1 MB RAM.
PowerPoint 3.0
For Windows, May 1992; for Mac: September 1992
Part of Microsoft Office for Windows 3.0 and Microsoft Office for Mac 3.0. Innovations included: the first application designed exclusively for the new Windows 3.1 platform, full support for TrueType fonts (new in Windows 3.1), presentation templates, editing in outline view, new drawing, including freeform tool, autoshapes, flip, rotate, scale, align, and transforming imported pictures into their drawing primitives to make them editable, transitions between slides in slide show, progressive builds, incorporating sound and video. Animations included "flying bullets" where bullet points "flew" into the slide one by one, and some degree of Pen Computing support was included.
"It added video-out to feed the new video projectors, with effects that could replace a bank of synchronized slide projectors. This version added fades, dissolves, and other transitions, as well as animation of text and pictures, and could incorporate video clips with synchronized audio."
System requirements: (Windows) 286 PC or higher, Windows 3.1, 2 MB RAM. (Mac) Macintosh Plus or better, System 7 or higher, 4 MB RAM.
PowerPoint 4.0
For Windows: February 1994; for Mac: October 1994
Part of Microsoft Office for Windows 4.0 and Microsoft Office for Mac 4.2. Innovations included: autolayouts, Word tables, rehearsal mode, hidden slides, and the "AutoContent Wizard".
Introduced a standard "Microsoft Office" look and feel (shared with Word and Excel), with status bar, toolbars, tooltips. Full OLE 2.0 with in-place activation.
System requirements: (Windows) 386 PC or higher, Windows 3.1, 8 MB RAM. (Mac) 68020 Mac or better, System 7 or higher, 8 MB RAM.
PowerPoint 7.0
For Windows: July 1995
Part of Microsoft Office for Windows 95. Innovations included: new animation effects, real curves and textures, black and white view, autocorrect, insert symbol, meeting support features such as "Meeting Minder".
"A complete rewrite of the product from the ground up in C++, full object model with internal VBA programmability".
System requirements: (Windows) 386 DX PC or higher, Windows 95, 6 MB RAM.
PowerPoint 8.0
For Windows: January 1997; for Mac: March 1998
Part of Microsoft Office for Windows 97 and Microsoft Office 98 Macintosh Edition. Innovations included: "Office Assistant", file compression, save to HTML, "Pack and Go", "AutoClipArt", transparent GIFs.
System requirements: (Windows) 486 PC or higher, 8 MB RAM. (Mac) PowerPC Mac or better, 16 MB RAM.
PowerPoint 9.0
For Windows: June 1999; for Mac: August 2000
Part of Microsoft Office for Windows 2000 and Microsoft Office for Mac 2001. Innovations included: three-pane "browser" view (selectable list of slide miniatures or titles, large single slide, notes), autofit text, real tables, presentation conferencing, save to web, picture bullets, animated GIFs, aliased fonts.
System requirements: (Windows) Pentium 75MHz+, Windows 95 or higher, 20 MB RAM. (Mac) PowerPC Mac 120MHz+ or better, MacOS 8.5 or higher, minimum 48 MB RAM.
PowerPoint 10.0
For Windows: May 2001; for Mac: November 2001
Part of Microsoft Office for Windows XP and Microsoft Office for Mac v.X. Innovations included: install from web, most clipart on web, use of Exchange and SharePoint for storage and collaboration.
System requirements: (Windows) Pentium III, Windows 98 or higher, 40 MB RAM. (Mac) OS X 10.1 ("Puma") or later (will not run under OS 9).
PowerPoint 11.0
For Windows: October 2003; for Mac: June 2004; for Mobile: May 2005
Part of Microsoft Office for Windows 2003 and Microsoft Office for Mac 2004. Innovations included: tools visible to presenter during slide show (notes, thumbnails, time clock, re-order and edit slides), "Package for CD" to write presentation and viewer app to CD. "Microsoft Producer for PowerPoint 2003" was a free plug-in from Microsoft, using a video camera, "that creates Web page presentations, with talking head narration, coordinated and timed to your existing PowerPoint presentation" for delivery over the web. The Genigraphics software to send a presentation for imaging as 35mm slides was removed from this version.
System requirements: (Windows) Pentium 233Mhz+, Windows 2000 with SP3 or later, 128 MB RAM. (Mac) Power Mac G3 or better, OS X 10.2.8 or later, 256 MB RAM.
PowerPoint 12.0
For Windows: January 2007; for Mobile: September 2007; for Mac: January 2008
Part of Microsoft Office for Windows 2007 and Microsoft Office for Mac 2008. Innovations included: new user interface ("Office Fluent") employing a changeable "ribbon" of tools across the top to replace menus and toolbars, SmartArt graphics, many graphical improvements in text and drawing, improved "Presenter View" (from 2003), widescreen slide formats. The "AutoContent Wizard" was removed from this version.
A major change in PowerPoint 2007 was from a binary file format, used from 1997 to 2003, to a new XML file format which evolved over further versions.
System requirements: (Windows) 500 MHz processor or higher, Windows XP with SP2 or later, 256 MB RAM. (Mac) 500 MHz processor or higher, MacOS X 10.4.9 or later, 512 MB RAM.
PowerPoint 14.0
For Windows: June 2010; for Web: June 2010; for Mobile: June 2010; for Mac: November 2010, for Symbian: April 2012
Part of Microsoft Office for Windows 2010 and Microsoft Office for Mac 2011. Innovations included: Single document interface (SDI), sections within presentations, reading view, redesign of "Backstage" functions (under File menu), save as video, insert video from web, embed video and audio, enhanced editing for video and for pictures, broadcast slideshow.
System requirements: (Windows) 500 MHz processor or higher, Windows XP with SP3 or later, 256 MB RAM, 512 MB RAM recommended for video. (Mac) Intel processor, Mac OS X 10.5.8 or later, 1 GB RAM.
PowerPoint 15.0
For Web: October 2012; for Mobile: November 2012; for Windows RT: November 2012; for Windows: January 2013; for iPhone: June 2013; for Android: July 2013; for Web: February 2014; for iPad: March 2014; for iOS: November 2014; for Mac: July 2015
Part of Microsoft Office for Windows 2013 and Microsoft Office for Mac 2016. Innovations included: Change default slide shape to 16:9 aspect ratio, online collaboration by multiple authors, user interface redesigned for multi-touch screens, improved audio, video, animations, and transitions, further changes to Presenter View. Clipart collections (and insertion tool) were removed, but available online.
System requirements: (Windows) 1 GHz processor or faster, x86- or x64-bit processor with SSE2 instruction set, Windows 7 or later, 1 GB RAM (32-bit), 2 GB RAM (64-bit). (Mac) Intel processor, Mac OS X 10.10 or later, 4 GB RAM.
PowerPoint 16.0
For Android: June 2015; for Mobile: July 2015; for iOS: July 2015; for Windows: September 2015; and Windows Store: January 2018
Part of Microsoft Office for Windows 2016. Innovations included: "Tell me" to search for program controls, "PowerPoint Designer" pane, Morph transition, real-time collaboration, "Zoom" to slides or sections in slideshow, and "Presentation Translator" for real-time translation of a presenter's spoken words to on-screen captions in any of 60+ languages, with the system analyzing the text of the PowerPoint presentation as context to increase the accuracy and relevance of the translations.
System requirements: (Windows) 1 GHz processor or faster, x86- or x64-bit processor with SSE2 instruction set, Windows 7 with SP 1 or later, 2 GB RAM.
File formats
Binary (1987–2007)
Early versions of PowerPoint, from 1987 through 1995 (versions 1.0 through 7.0), evolved through a sequence of binary file formats, different in each version, as functionality was added. This set of formats were never documented, but an open-source libmwaw (used by LibreOffice) exists to read them.
A stable binary format (called a .ppt file, like all earlier binary formats) that was shared as the default in PowerPoint 97 through PowerPoint 2003 for Windows, and in PowerPoint 98 through PowerPoint 2004 for Mac (that is, in PowerPoint versions 8.0 through 11.0) was finally created. It was based on the Compound File Binary Format. The specification document is actively maintained and can be freely downloaded, because, although no longer the default, that binary format can be read and written by some later versions of PowerPoint, including PowerPoint 2016. After the stable binary format was adopted, versions of PowerPoint continued to be able to read and write differing file formats from earlier versions. But beginning with PowerPoint 2007 and PowerPoint 2008 for Mac (PowerPoint version 12.0), this was the only binary format available for saving; PowerPoint 2007 (version 12.0) no longer supported saving to binary file formats used earlier than PowerPoint 97 (version 8.0), ten years before.
The ".pps" and ".ppsx" file extensions are technically the same as ".ppt" and ".pptx", except they are launched as presentation instead of for editing by default.
Binary filename extensions
.ppt, PowerPoint 97–2003 binary presentation
.pps, PowerPoint 97–2003 binary slide show
.pot, PowerPoint 97–2003 binary template
Binary media types
.ppt, application/vnd.ms-powerpoint
.pps, application/vnd.ms-powerpoint
.pot, application/vnd.ms-powerpoint
Office Open XML (since 2007)
The big change in PowerPoint 2007 and PowerPoint 2008 for Mac (PowerPoint version 12.0) was that the stable binary file format of 97–2003 was replaced as the default by a new zipped XML-based Office Open XML format (.pptx files). Microsoft's explanation of the benefits of the change included: smaller file sizes, up to 75% smaller than comparable binary documents; security, through being able to identify and exclude executable macros and personal data; less chance to be corrupted than binary formats; and easier interoperability for exchanging data among Microsoft and other business applications, all while maintaining backward compatibility.
XML filename extensions
.pptx, PowerPoint 2007 XML presentation
.pptm, PowerPoint 2007 XML macro-enabled presentation
.ppsx, PowerPoint 2007 XML slide show
.ppsm, PowerPoint 2007 XML macro-enabled slide show
.ppam, PowerPoint 2007 XML add-in
.potx, PowerPoint 2007 XML template
.potm, PowerPoint 2007 XML macro-enabled template
XML media types
.pptx, application/vnd.openxmlformats-officedocument.presentationml.presentation
.pptm, application/vnd.ms-powerpoint.presentation.macroEnabled.12
.ppsx, application/vnd.openxmlformats-officedocument.presentationml.slideshow
.ppsm, application/vnd.ms-powerpoint.slideshow.macroEnabled.12
.ppam, application/vnd.ms-powerpoint.addin.macroEnabled.12
.potx, application/vnd.openxmlformats-officedocument.presentationml.template
.potm, application/vnd.ms-powerpoint.template.macroEnabled.12
The specification for the new format was published as an open standard, ECMA-376, through Ecma International Technical Committee 45 (TC45). The Ecma 376 standard was approved in December 2006, and was submitted for standardization through ISO/IEC JTC 1/SC 34 WG4 in early 2007. The standardization process was contentious. It was approved as ISO/IEC 29500 in early 2008. Copies of the ISO/IEC standard specification are freely available, in two parts. These define two related standards known as "Transitional" and "Strict". The two standards were progressively adopted by PowerPoint: PowerPoint version 12.0 (2007, 2008 for Mac) could read and write Transitional format, but could neither read nor write Strict format. PowerPoint version 14.0 (2010, 2011 for Mac) could read and write Transitional, and also read but not write Strict. PowerPoint version 15.0 and later (beginning 2013, 2016 for Mac) can read and write both Transitional and Strict formats. The reason for the two variants was explained by Microsoft:
The PowerPoint .pptx file format (called "PresentationML" for Presentation Markup Language) contains separate structures for all the complex parts of a PowerPoint presentation. The specification documents run to over six thousand pages. Because of the widespread use of PowerPoint, the standardized file formats are considered important for the long-term access to digital documents in library collections and archives, according to the U.S. Library of Congress.
PowerPoint 2013 and PowerPoint 2016 provide options to set default saving to ISO/IEC 29500 Strict format, but the initial default setting remains Transitional, for compatibility with legacy features incorporating binary data in existing documents. PowerPoint 2013 or PowerPoint 2016 will both open and save files in the former binary format (.ppt), for compatibility with older versions of the program (but not versions older than PowerPoint 97). In saving to older formats, these versions of PowerPoint will check to assure that no features have been introduced into the presentation which are incompatible with the older formats.
PowerPoint 2013 and 2016 will also save a presentation in many other file formats, including PDF format, MPEG-4 or WMV video, as a sequence of single-picture files (using image formats including GIF, JPEG, PNG, TIFF, and some older formats), and as a single presentation file in which all slides are replaced with pictures. PowerPoint will both open and save files in OpenDocument Presentation format (ODP) for compatibility.
| Technology | Office and data management | null |
77288 | https://en.wikipedia.org/wiki/Common%20nightingale | Common nightingale | The common nightingale, rufous nightingale or simply nightingale (Luscinia megarhynchos), is a small passerine bird which is best known for its powerful and beautiful song. It was formerly classed as a member of the thrush family Turdidae, but is now more generally considered to be an Old World flycatcher, Muscicapidae. It belongs to a group of more terrestrial species, often called chats.
Etymology
"Nightingale" is derived from "night" and the Old English galan, "to sing". The genus name Luscinia is Latin for "nightingale" and megarhynchos is from Ancient Greek megas, "great" and rhunkhos "bill".
Subspecies
western nightingale (L. m. megarhynchos) - Western Europe, North Africa and Asia Minor, wintering in tropical Africa
Caucasian nightingale (L. m. africana) - The Caucasus and eastern Turkey to southwestern Iran and Iraq, wintering in East Africa
eastern nightingale (L. m. golzii) - The Aral Sea to Mongolia, wintering in coastal East Africa
Description
The common nightingale is slightly larger than the European robin, at length. It is plain brown above except for the reddish tail. It is buff to white below. The sexes are similar. The eastern subspecies (L. m. golzi) and the Caucasian subspecies (L. m. africana) have paler upper parts and a stronger face-pattern, including a pale supercilium. The song of the male nightingale has been described as one of the most beautiful sounds in nature, inspiring songs, fairy tales, opera, books, and a great deal of poetry. However, historically most people were not aware that female nightingales do not sing.
Distribution and habitat
It is a migratory insectivorous species breeding in forest and scrub in Europe and the Palearctic, and wintering in Sub-Saharan Africa. It is not found naturally in the Americas. The distribution is more southerly than the very closely related thrush nightingale Luscinia luscinia. It nests on or near the ground in dense vegetation. Research in Germany found that favoured breeding habitat of nightingales was defined by a number of geographical factors.
less than above mean sea level
mean air temperature during the growing season above
more than 20 days/year on which temperatures exceed
annual precipitation less than
aridity index lower than 0.35
no closed canopy
In the U.K., the bird is at the northern limit of its range which has contracted in recent years, placing it on the red list for conservation. Despite local efforts to safeguard its favoured coppice and scrub habitat, numbers fell by 53 percent between 1995 and 2008. A survey conducted by the British Trust for Ornithology in 2012 and 2013 recorded some 3,300 territories, with most of these clustered in a few counties in the southeast of England, notably Kent, Essex, Suffolk, and East and West Sussex.
By contrast, the European breeding population is estimated at between 3.2 and 7 million pairs, giving it green conservation status (least concern).
Behaviour and ecology
Common nightingales are so named because they frequently sing at night as well as during the day. The name has been used for more than 1,000 years, being highly recognisable even in its Old English form nihtegale, which means "night songstress". Early writers assumed the female sang when it is in fact the male. The song is loud, with an impressive range of whistles, trills and gurgles. Its song is particularly noticeable at night because few other birds are singing. This is why its name includes "night" in several languages. Only unpaired males sing regularly at night, and nocturnal song probably serves to attract a mate. Singing at dawn, during the hour before sunrise, is assumed to be important in defending the bird's territory. Nightingales sing even more loudly in urban or near-urban environments, in order to overcome the background noise. The most characteristic feature of the song is a loud whistling crescendo that is absent from the song of its close relative, the thrush nightingale (Luscinia luscinia). It has a frog-like alarm call.
The bird is a host of the acanthocephalan intestinal parasite Apororhynchus silesiacus.
Cultural connotations
The common nightingale is an important symbol for poets from a variety of ages, and has taken on a number of symbolic connotations. Homer evokes Aëdon the nightingale in Odyssey, suggesting the myth of Philomela and Procne (one of whom, depending on the myth's version, is turned into a nightingale). This myth is the focus of Sophocles' tragedy, Tereus, of which only fragments remain. Ovid, too, in his Metamorphoses, includes the most popular version of this myth, imitated and altered by later poets, including Chrétien de Troyes, Geoffrey Chaucer, John Gower, and George Gascoigne. T.S. Eliot's "The Waste Land" also evokes the common nightingale's song (and the myth of Philomela and Procne). Because of the violence associated with the myth, the nightingale's song was long interpreted as a lament.
The common nightingale has also been used as a symbol of poets or their poetry. Poets chose the nightingale as a symbol because of its creative and seemingly spontaneous song. Aristophanes's The Birds and Callimachus both evoke the bird's song as a form of poetry. Virgil compares the mourning of Orpheus to the “lament of the nightingale”.
In Sonnet 102 Shakespeare compares his love poetry to the song of the common nightingale (Philomel):
"Our love was new, and then but in the spring,
When I was wont to greet it with my lays;
As Philomel in summer's front doth sing,
And stops his pipe in growth of riper days:"
During the Romantic era the bird's symbolism changed once more: poets viewed the nightingale not only as a poet in his own right, but as “master of a superior art that could inspire the human poet”. For some romantic poets, the nightingale even began to take on qualities of the muse. The nightingale has a long history with symbolic associations ranging from "creativity, the muse, nature's purity, and, in Western spiritual tradition, virtue and goodness." Coleridge and Wordsworth saw the nightingale more as an instance of natural poetic creation: the nightingale became a voice of nature. John Keats' "Ode to a Nightingale" pictures the nightingale as an idealized poet who has achieved the poetry that Keats longs to write. Invoking a similar conception of the nightingale, Shelley wrote in his “A Defence of Poetry":
A poet is a nightingale who sits in darkness and sings to cheer its own solitude with sweet sounds; his auditors are as men entranced by the melody of an unseen musician, who feel that they are moved and softened, yet know not whence or why.
The nightingale is the national bird of Ukraine. One legend tells how nightingales once only lived in India, when one nightingale visited Ukraine. Hearing sad songs from the people, the nightingale sang its song to cheer them up. The people responded with happy songs, and since then, nightingales have visited Ukraine every spring to hear Ukrainian songs. National poet Taras Shevchenko observed that "even the memory of the nightingale's song makes man happy."
The nightingale is the official national bird of Iran. In medieval Persian literature, the nightingale's enjoyable song has made it a symbol of the lover who is eloquent, passionate, and doomed to love in vain. In Persian poetry, the object of the nightingale's affections is the rose which embodies both the perfection of earthly beauty and the arrogance of that perfection.
Cultural depictions
The Aēdōn (, "Nightingale") is a minor character in Aristophanes's 414 BC Attic comedy The Birds.
Philomela is transformed into a nightingale, according to Metamorphoses (book VI) of Ovid.
The love of the nightingale (a conventional cultural substitution for the Persian bulbul) for the rose is widely used as a metaphor for the poet's love for the beloved and the worshiper's love for God in classical Persian, Urdu and Turkish poetry.
"The Owl and the Nightingale" (12th or 13th century) is a Middle English poem about an argument between these two birds.
"When The Nightingale Sings" is a Middle English love poem, extolling the beauty and lost love of an unknown maiden.
"Laüstic", a lai by French poet Marie de France from High Middle Ages (1100–1300)
John Milton's sonnet "To the Nightingale" (1632–33) contrasts the symbolism of the nightingale as a bird for lovers, with the cuckoo as the bird that called when wives were unfaithful to (or "cuckolded") their husbands.
Samuel Taylor Coleridge's "The Nightingale: A Conversation Poem", printed in 1798, disputes the traditional idea that nightingales are connected to the idea of melancholy.
Ludwig van Beethoven's Symphony No. 6 (1808), the "Pastoral Symphony", includes in its second movement flute imitations of nightingale calls.
A nightingale (called by its French name rossignol) features prominently in the French folk song À la claire fontaine
Franz Liszt featured the nightingale's song in the Mephisto Waltzes No. 1.
John Keats' "Ode to a Nightingale" (1819) was described by Edmund Clarence Stedman as "one of our shorter English lyrics that still seems to me... the nearest to perfection, the one I would surrender last of all" and by Algernon Charles Swinburne as "one of the final masterpieces of human work in all time and for all ages".
The beauty of the nightingale's song is a theme in Hans Christian Andersen's story "The Nightingale" from 1843.
A recording of nightingale song is included, as directed by the score, in "The Pines of Janiculum", the third movement of Ottorino Respighi's 1924 symphonic poem Pines of Rome ().
Igor Stravinsky based his first opera, The Nightingale (1914), on the Hans Christian Andersen story and later prepared a symphonic poem, The Song of the Nightingale (1917), using music from the opera.
In 1915, Joseph Lamb wrote a rag called "Ragtime Nightingale" that was intended to imitate the nightingale calls.
"A Nightingale Sang in Berkeley Square" (1939) was one of the most popular songs in Britain during World War II. In 2004, the song was featured in an episode of series 2 of the Channel 4 sitcom Peep Show and in 2019, it featured as the closing song of the Amazon/BBC miniseries Good Omens.
Both Terry Pratchett and Neil Gaiman's novel Good Omens and the aforementioned miniseries adaptation joke that, "...while they were eating, for the first time ever, a nightingale (sang/actually did sing) in Berkeley Square. Nobody heard it over the noise of the traffic, but it was there, right enough."
In the works of J. R. R. Tolkien, nightingales are closely associated with the characters Lúthien Tinúviel and her mother, Melian.
A nightingale is depicted on the reverse of the Croatian 1 kuna coin, minted between 1993 and 2009.
Nightingale was an inspiration of the creation of a Korean court solo dance Chunaengjeon (춘앵전). The dance initially was performed by a female dancer of the court of Joseon Dynasty, Mudong.
In Chapter 13 of Mary Shelley's Frankenstein, the monster compares Safie's singing voice to that of a "nightingale in the woods".
Manfred Mann's Earth Band's sixth album, 1975's Nightingales & Bombers, took its title from a World War II naturalist's recording of a nightingale singing in a garden as warplanes flew overhead. The recording is featured in a song on the album.
The song of the nightingale is one of the main elements in the 2019 single "Let Nature Sing".
An operator in the mobile video game Arknights is named after it.
In the Baha'i Faith
The nightingale is used symbolically in the Baha'i Faith to represent the founder Baha'u'llah. Baha'is utilise this metaphor to convey how Baha'u'llah's writings are of beautiful quality, much like how the nightingale's singing is revered for its beautiful quality in Persian music and literature.
Nightingales are mentioned in much of Baha'u'llah's works, including the Tablet of Ahmad, The Seven Valleys, The Hidden Words, and the untranslated Tablet of the Nightingale and the Owl.
| Biology and health sciences | Passerida | null |
77305 | https://en.wikipedia.org/wiki/Flea | Flea | Flea, the common name for the order Siphonaptera, includes 2,500 species of small flightless insects that live as external parasites of mammals and birds. Fleas live by ingesting the blood of their hosts. Adult fleas grow to about long, are usually brown, and have bodies that are "flattened" sideways or narrow, enabling them to move through their hosts' fur or feathers. They lack wings; their hind legs are extremely well adapted for jumping. Their claws keep them from being dislodged, and their mouthparts are adapted for piercing skin and sucking blood. Some species can leap 50 times their body length, a feat second only to jumps made by another group of insects, the superfamily of froghoppers. Flea larvae are worm-like, with no limbs; they have chewing mouthparts and feed on organic debris left on their hosts' skin.
Genetic evidence indicates that fleas are a specialised lineage of parasitic scorpionflies (Mecoptera) sensu lato, most closely related to the family Nannochoristidae. The earliest known fleas lived in the Middle Jurassic; modern-looking forms appeared in the Cenozoic. Fleas probably originated on mammals first and expanded their reach to birds. Each species of flea specializes, more or less, on one species of host: many species of flea never breed on any other host; some are less selective. Some families of fleas are exclusive to a single host group; for example, the Malacopsyllidae are found only on armadillos, the Ischnopsyllidae only on bats, and the Chimaeropsyllidae only on elephant shrews.
The oriental rat flea, Xenopsylla cheopis, is a vector of Yersinia pestis, the bacterium that causes bubonic plague. The disease was spread to humans by rodents, such as the black rat, which were bitten by infected fleas. Major outbreaks included the Plague of Justinian, about 540, and the Black Death, about 1350, each of which killed a sizeable fraction of the world's people.
Fleas appear in human culture in such diverse forms as flea circuses; poems, such as John Donne's erotic "The Flea"; works of music, such as those by Modest Mussorgsky; and a film by Charlie Chaplin.
Morphology and behavior
Fleas are wingless insects, long, that are agile, usually dark colored (for example, the reddish-brown of the cat flea), with a proboscis, or stylet, adapted to feeding by piercing the skin and sucking their host's blood through their epipharynx. Flea legs end in strong claws that are adapted to grasp a host.
Unlike other insects, fleas do not possess compound eyes but instead only have simple eyespots with a single biconvex lens; some species lack eyes altogether. Their bodies are laterally compressed, permitting easy movement through the hairs or feathers on the host's body. The flea body is covered with hard plates called sclerites. These sclerites are covered with many hairs and short spines directed backward, which also assist its movements on the host. The tough body is able to withstand great pressure, likely an adaptation to survive attempts to eliminate them by scratching.
Fleas lay tiny, white, oval eggs. The larvae are small and pale, have bristles covering their worm-like bodies, lack eyes, and have mouth parts adapted to chewing. The larvae feed on organic matter, especially the feces of mature fleas, which contain dried blood. Adults feed only on fresh blood.
Jumping
Their legs are long, the hind pair well adapted for jumping; a flea can jump vertically up to and horizontally up to , making the flea one of the best jumpers of all known animals (relative to body size), second only to the froghopper. A flea can jump 60 times its length in height and 110 times its length in distance, equivalent to a adult human jumping vertically and horizontally. Rarely do fleas jump from dog to dog. Most flea infestations come from newly developed fleas from the pet's environment. The flea jump is so rapid and forceful that it exceeds the capabilities of muscle, and instead of relying on direct muscle power, fleas store muscle energy in a pad of the elastic protein named resilin before releasing it rapidly (like a human using a bow and arrow). Immediately before the jump, muscles contract and deform the resilin pad, slowly storing energy which can then be released extremely rapidly to power leg extension for propulsion. To prevent premature release of energy or motions of the leg, the flea employs a "catch mechanism". Early in the jump, the tendon of the primary jumping muscle passes slightly behind the coxa-trochanter joint, generating a torque which holds the joint closed with the leg close to the body. To trigger jumping, another muscle pulls the tendon forward until it passes the joint axis, generating the opposite torque to extend the leg and power the jump by release of stored energy. The actual take off has been shown by high-speed video to be from the tibiae and tarsi rather than from the trochantera (knees).
Life cycle and development
Fleas are holometabolous insects, going through the four lifecycle stages of egg, larva, pupa, and imago (adult). In most species, neither female nor male fleas are fully mature when they first emerge but must feed on blood before they become capable of reproduction. The first blood meal triggers the maturation of the ovaries in females and the dissolution of the testicular plug in males, and copulation soon follows. Some species breed all year round while others synchronise their activities with their hosts' life cycles or with local environmental factors and climatic conditions. Flea populations consist of roughly 50% eggs, 35% larvae, 10% pupae, and 5% adults.
Egg
The number of eggs laid depends on species, with batch sizes ranging from two to several dozen. The total number of eggs produced in a female's lifetime (fecundity) varies from around one hundred to several thousand. In some species, the flea lives in the host's nest or burrow and the eggs are deposited on the substrate, but in others, the eggs are laid on the host itself and can easily fall off onto the ground. Because of this, areas where the host rests and sleeps become one of the primary habitats of eggs and developing larvae. The eggs take around two days to two weeks to hatch.
Larva
Flea larvae emerge from the eggs to feed on any available organic material such as dead insects, faeces, conspecific eggs, and vegetable matter. In laboratory studies, some dietary diversity seems necessary for proper larval development. Blood-only diets allow only 12% of larvae to mature, whereas blood and yeast or dog chow diets allow almost all larvae to mature. Another study also showed that 90% of larvae matured into adults when the diet included nonviable eggs. They are blind and avoid sunlight, keeping to dark, humid places such as sand or soil, cracks and crevices, under carpets and in bedding. The entire larval stage lasts between four and 18 days.
Pupa
Given an adequate supply of food, larvae pupate and weave silken cocoons after three larval stages. Within the cocoon, the larva molts for a final time and undergoes metamorphosis into the adult form. This can take just four days, but may take much longer under adverse conditions, and there follows a variable-length stage during which the pre-emergent adult awaits a suitable opportunity to emerge. Trigger factors for emergence include vibrations (including sound), heat (in warm-blooded hosts), and increased levels of carbon dioxide, all of which may indicate the presence of a suitable host. Large numbers of pre-emergent fleas may be present in otherwise flea-free environments, and the introduction of a suitable host may trigger a mass emergence.
Adult
Once the flea reaches adulthood, its primary goal is to find blood and then to reproduce. Female fleas can lay 5000 or more eggs over their life, permitting rapid increase in numbers. Generally speaking, an adult flea only lives for 2 or 3 months. Without a host to provide a blood meal, a flea's life can be as short as a few days. Under ideal conditions of temperature, food supply, and humidity, adult fleas can live for up to a year and a half. Completely developed adult fleas can live for several months without eating, so long as they do not emerge from their puparia. Optimum temperatures for the flea's life cycle are and optimum humidity is 70%.
Adult female rabbit fleas, Spilopsyllus cuniculi, can detect the changing levels of cortisol and corticosterone hormones in the rabbit's blood that indicate it is getting close to giving birth. This triggers sexual maturity in the fleas and they start producing eggs. As soon as the baby rabbits are born, the fleas make their way down to them and once on board they start feeding, mating, and laying eggs. After 12 days, the adult fleas make their way back to the mother. They complete this mini-migration every time she gives birth.
Taxonomy and phylogeny
History
Between 1735 and 1758, the Swedish naturalist Carl Linnaeus first classified insects, doing so on the basis of their wing structure. One of the seven orders into which he divided them was "Aptera", meaning wingless, a group in which as well as fleas, he included spiders, woodlice and myriapods. It wasn't until 1810 that the French zoologist Pierre André Latreille reclassified the insects on the basis of their mouthparts as well as their wings, splitting Aptera into Thysanura (silverfish), Anoplura (sucking lice) and Siphonaptera (fleas), at the same time separating off the arachnids and crustaceans into their own subphyla. The group's name, Siphonaptera, is zoological Latin from the Greek siphon (a tube) and aptera (wingless).
External phylogeny
It was historically unclear whether the Siphonaptera are sister to the Mecoptera (scorpionflies and allies), or are inside that clade, making "Mecoptera" paraphyletic. The earlier suggestion that the Siphonaptera are sister to the Boreidae (snow scorpionflies) is not supported. A 2020 genetic study recovered Siphonaptera within Mecoptera, with strong support, as the sister group to Nannochoristidae, a small, relictual group of mecopterans native to the Southern Hemisphere. Fleas and nannochoristids share several similarities with each other that are not shared with other mecopterans, including similar mouthparts as well as a similar sperm pump organisation.
Relationships of Siphonaptera per Tihelka et al. 2020.
Fossil history
Fleas likely descended from scorpionflies, insects that are predators or scavengers. Fossils of large, wingless stem-group fleas with siphonate (sucking) mouthparts from the Middle Jurassic to Early Cretaceous have been found in northeastern China and Russia, belonging to the families Saurophthiridae and Pseudopulicidae, as well as Tarwinia from the Early Cretaceous of Australia. Most flea families formed after the end of the Cretaceous (in the Paleogene and onwards). Modern fleas probably arose in the southern continental area of Gondwana, and migrated rapidly northwards from there. They most likely evolved with mammal hosts, only later moving to birds.
Siphonaptera is a relatively small order of insects: members of the order undergo complete metamorphosis and are secondarily wingless (their ancestors had wings which modern forms have lost). In 2005, Medvedev listed 2005 species in 242 genera, and despite subsequent descriptions of new species, bringing the total up to around 2500 species, this is the most complete database available. The order is divided into four infraorders and eighteen families. Some families are exclusive to a single host group; these include the Malacopsyllidae (armadillos), Ischnopsyllidae (bats) and Chimaeropsyllidae (elephant shrews).
Many of the known species are little studied. Some 600 species (a quarter of the total) are known from single records. Over 94% of species are associated with mammalian hosts, and only about 3% of species can be considered to be specific parasites of birds. The fleas on birds are thought to have originated from mammalian fleas; at least sixteen separate groups of fleas switched to avian hosts during the evolutionary history of the Siphonaptera. Occurrences of fleas on reptiles is accidental, and fleas have been known to feed on the hemolymph (bloodlike body fluid) of ticks.
Internal phylogeny
Flea phylogeny was long neglected, the discovery of homologies with the parts of other insects being made difficult by their extreme specialization. Whiting and colleagues prepared a detailed molecular phylogeny in 2008, with the basic structure shown in the cladogram. The Hectopsyllidae, including the harmful chigoe flea or jigger, is sister to the rest of the Siphonaptera.
Taxonomy
, there are 21 recognized families within the order Siphonaptera, 3 of which are extinct. In addition, some researchers have suggested that the subfamily Stenoponiinae should be elevated to its own family (Stenoponiidae).
Ancistropsyllidae Toumanoff & Fuller, 1947
Ceratophyllidae Dampf, 1908
Chimaeropsyllidae Ewing & I. Fox, 1943
Coptopsyllidae Wagner, 1928
Ctenophthalmidae Rothschild, 1915
Hystrichopsyllidae Tiraboschi, 1904
Ischnopsyllidae Wahlgren, 1907
Leptopsyllidae Rothschild & Jordan, 1915
Lycopsyllidae Baker, 1905
Malacopsyllidae Baker, 1905
Pseudopulicidae† Gao, Shih & Ren, 2012
Pulicidae Billberg, 1820
Pygiopsyllidae Wagner, 1939
Rhopalopsyllidae Oudemans, 1909
Saurophthiridae† Ponomarenko, 1986
Stephanocircidae Wagner, 1928
Stivaliidae Mardon, 1978
Tarwiniidae† Huang, Engel, Cai & Nel, 2013
Tungidae Fox, 1925
Vermipsyllidae Wagner, 1889
Xiphiopsyllidae Wagner, 1939
Relationship with host
Fleas feed on a wide variety of warm-blooded vertebrates including dogs, cats, rabbits, squirrels, ferrets, rats, mice, birds, and sometimes humans. Fleas normally specialise in one host species or group of species, but can often feed but not reproduce on other species. Ceratophyllus gallinae affects poultry as well as wild birds. As well as the degree of relatedness of a potential host to the flea's original host, it has been shown that avian fleas that exploit a range of hosts, only parasitise species with low immune responses. In general, host specificity decreases as the size of the host species decreases. Another factor is the opportunities available to the flea to change host species; this is smaller in colonially nesting birds, where the flea may never encounter another species, than it is in solitary nesting birds. A large, long-lived host provides a stable environment that favours host-specific parasites.
Although there are species named dog fleas (Ctenocephalides canis Curtis, 1826) and cat fleas (Ctenocephalides felis), fleas are not always strictly species-specific. A study in Virginia examined 244 fleas from 29 dogs: all were cat fleas. Dog fleas had not been found in Virginia in more than 70 years, and may not even occur in the US, so a flea found on a dog is likely a cat flea (Ctenocephalides felis).
One theory of human hairlessness is that the loss of hair helped humans to reduce their burden of fleas and other ectoparasites.
Direct effects of bites
In many species, fleas are principally a nuisance to their hosts, causing an itching sensation which in turn causes the host to try to remove the pest by biting, pecking or scratching. Fleas are not simply a source of annoyance, however. Flea bites cause a slightly raised, swollen, irritating nodule to form on the epidermis at the site of each bite, with a single puncture point at the centre, like a mosquito bite. This can lead to an eczematous itchy skin disease called flea allergy dermatitis, which is common in many host species, including dogs and cats. The bites often appear in clusters or lines of two bites, and can remain itchy and inflamed for up to several weeks afterwards. Fleas can lead to secondary hair loss as a result of frequent scratching and biting by the animal. They can also cause anemia in extreme cases.
As a vector
Fleas are vectors for viral, bacterial and rickettsial diseases of humans and other animals, as well as of protozoan and helminth parasites. Bacterial diseases carried by fleas include murine or endemic typhus and bubonic plague. Fleas can transmit Rickettsia typhi, Rickettsia felis, Bartonella henselae, and the myxomatosis virus. They can carry Hymenolepiasis tapeworms and Trypanosome protozoans. The chigoe flea or jigger (Tunga penetrans) causes the disease tungiasis, a major public health problem around the world. Fleas that specialize as parasites on specific mammals may use other mammals as hosts; thus, humans may be bitten by cat and dog fleas.
Relationship with humans
In literature and art
Fleas have appeared in poetry, literature, music and art; these include Robert Hooke's drawing of a flea under the microscope in his pioneering book Micrographia published in 1665, poems by Donne and Jonathan Swift, works of music by Giorgio Federico Ghedini and Modest Mussorgsky, a play by Georges Feydeau, a film by Charlie Chaplin, and paintings by artists such as Giuseppe Crespi, Giovanni Battista Piazzetta, and Georges de La Tour.
John Donne's erotic metaphysical poem "The Flea", published in 1633 after his death, uses the conceit of a flea, which has sucked blood from the male speaker and his female lover, as an extended metaphor for their sexual relationship. The speaker tries to convince a lady to sleep with him, arguing that if the mingling of their blood in the flea is innocent, then sex would be also.
The comic poem Siphonaptera was written in 1915 by the mathematician Augustus De Morgan, It describes an infinite chain of parasitism made of ever larger and ever smaller fleas.
Flea circuses
Flea circuses provided entertainment to nineteenth century audiences. These circuses, extremely popular in Europe from 1830 onwards, featured fleas dressed as humans or towing miniature carts, chariots, rollers or cannon. These devices were originally made by watchmakers or jewellers to show off their skill at miniaturization. A ringmaster called a "professor" accompanied their performance with a rapid circus patter.
Carriers of plague
Oriental rat fleas, Xenopsylla cheopis, can carry the coccobacillus Yersinia pestis. The infected fleas feed on rodent vectors of this bacterium, such as the black rat, Rattus rattus, and then infect human populations with the plague, as has happened repeatedly from ancient times, as in the Plague of Justinian in 541–542. Outbreaks killed up to 200 million people across Europe between 1346 and 1671. The Black Death pandemic between 1346 and 1353 likely killed over a third of the population of Europe.
Because fleas carry plague, they have seen service as a biological weapon. During World War II, the Japanese army dropped fleas infested with Y. pestis in China. The bubonic and septicaemic plagues are the most probable form of the plague that would spread as a result of a bioterrorism attack that used fleas as a vector.
The Rothschild Collection
The banker Charles Rothschild devoted much of his time to entomology, creating a large collection of fleas now in the Rothschild Collection at the Natural History Museum, London. He discovered and named the plague vector flea, Xenopsylla cheopis, also known as the oriental rat flea, in 1903. Using what was probably the world's most complete collection of fleas of about 260,000 specimens (representing some 73% of the 2,587 species and subspecies so far described), he described around 500 species and subspecies of Siphonaptera. He was followed in this interest by his daughter Miriam Rothschild, who helped to catalogue his enormous collection of the insects in seven volumes.
Flea treatments
Fleas have a significant economic impact. In America alone, approximately $2.8 billion is spent annually on flea-related veterinary bills and another $1.6 billion annually for flea treatment with pet groomers. Four billion dollars is spent annually for prescription flea treatment and $348 million for flea pest control.
| Biology and health sciences | Insects and other hexapods | null |
77385 | https://en.wikipedia.org/wiki/Polyethylene | Polyethylene | Polyethylene or polythene (abbreviated PE; IUPAC name polyethene or poly(methylene)) is the most commonly produced plastic. It is a polymer, primarily used for packaging (plastic bags, plastic films, geomembranes and containers including bottles, cups, jars, etc.). , over 100 million tonnes of polyethylene resins are being produced annually, accounting for 34% of the total plastics market.
Many kinds of polyethylene are known, with most having the chemical formula (C2H4)n. PE is usually a mixture of similar polymers of ethylene, with various values of n. It can be low-density or high-density and many variations thereof. Its properties can be modified further by crosslinking or copolymerization. All forms are nontoxic as well as chemically resilient, contributing to polyethylene's popularity as a multi-use plastic. However, polyethylene's chemical resilience also makes it a long-lived and decomposition-resistant pollutant when disposed of improperly. Being a hydrocarbon, polyethylene is colorless to opaque (without impurities or colorants) and combustible.
History
Polyethylene was first synthesized by the German chemist Hans von Pechmann, who prepared it by accident in 1898 while investigating diazomethane. When his colleagues Eugen Bamberger and Friedrich Tschirner characterized the white, waxy substance that he had created, they recognized that it contained long −CH2− chains and termed it polymethylene.
The first industrially practical polyethylene synthesis (diazomethane is a notoriously unstable substance that is generally avoided in industrial syntheses) was again accidentally discovered in 1933 by Eric Fawcett and Reginald Gibson at the Imperial Chemical Industries (ICI) works in Northwich, England. Upon applying extremely high pressure (several hundred atmospheres) to a mixture of ethylene and benzaldehyde they again produced a white, waxy material. Because the reaction had been initiated by trace oxygen contamination in their apparatus, the experiment was difficult to reproduce at first. It was not until 1935 that another ICI chemist, Michael Perrin, developed this accident into a reproducible high-pressure synthesis for polyethylene that became the basis for industrial low-density polyethylene (LDPE) production beginning in 1939. Because polyethylene was found to have very low-loss properties at very high frequency radio waves, commercial distribution in Britain was suspended on the outbreak of World War II, secrecy imposed, and the new process was used to produce insulation for UHF and SHF coaxial cables of radar sets. During World War II, further research was done on the ICI process and in 1944, DuPont at Sabine River, Texas, and Union Carbide Corporation at South Charleston, West Virginia, began large-scale commercial production under license from ICI.
The landmark breakthrough in the commercial production of polyethylene began with the development of catalysts that promoted the polymerization at mild temperatures and pressures. The first of these was a catalyst based on chromium trioxide discovered in 1951 by Robert Banks and J. Paul Hogan at Phillips Petroleum. In 1953 the German chemist Karl Ziegler developed a catalytic system based on titanium halides and organoaluminium compounds that worked at even milder conditions than the Phillips catalyst. The Phillips catalyst is less expensive and easier to work with, however, and both methods are heavily used industrially. By the end of the 1950s both the Phillips- and Ziegler-type catalysts were being used for high-density polyethylene (HDPE) production. In the 1970s, the Ziegler system was improved by the incorporation of magnesium chloride. Catalytic systems based on soluble catalysts, the metallocenes, were reported in 1976 by Walter Kaminsky and Hansjörg Sinn. The Ziegler- and metallocene-based catalysts families have proven to be very flexible at copolymerizing ethylene with other olefins and have become the basis for the wide range of polyethylene resins available today, including very-low-density polyethylene and linear low-density polyethylene. Such resins, in the form of UHMWPE fibers, have (as of 2005) begun to replace aramids in many high-strength applications.
Properties
The properties of polyethylene depend strongly on type. The molecular weight, crosslinking, and presence of comonomers all strongly affect its properties. It is for this structure-property relation that intense effort has been invested into diverse kinds of PE. LDPE is softer and more transparent than HDPE. For medium- and high-density polyethylene the melting point is typically in the range . The melting point for average commercial low-density polyethylene is typically . These temperatures vary strongly with the type of polyethylene, but the theoretical upper limit of melting of polyethylene is reported to be . Combustion typically occurs above .
Most LDPE, MDPE, and HDPE grades have excellent chemical resistance, meaning that they are not attacked by strong acids or strong bases and are resistant to gentle oxidants and reducing agents. Crystalline samples do not dissolve at room temperature. Polyethylene (other than cross-linked polyethylene) usually can be dissolved at elevated temperatures in aromatic hydrocarbons such as toluene or xylene, or in chlorinated solvents such as trichloroethane or trichlorobenzene.
Polyethylene absorbs almost no water; the gas and water vapour permeability (only polar gases) is lower than for most plastics. Oxygen, carbon dioxide and flavorings, on the other hand, can pass it easily.
Polyethylene burns slowly with a blue flame having a yellow tip and gives off an odour of paraffin (similar to candle flame). The material continues burning on removal of the flame source and produces a drip.
Polyethylene cannot be imprinted or bonded with adhesives without pretreatment. High-strength joints are readily achieved with plastic welding.
Electrical
Polyethylene is a good electrical insulator. It offers good electrical treeing resistance; however, it becomes easily electrostatically charged (which can be reduced by additions of graphite, carbon black or antistatic agents).
When pure, the dielectric constant is in the range 2.2 to 2.4 depending on the density and the loss tangent is very low, making it a good dielectric for building capacitors. For the same reason it is commonly used as the insulation material for high-frequency coaxial and twisted pair cables.
Optical
Depending on thermal history and film thickness, PE can vary between almost clear (transparent), milky-opaque (translucent) and opaque. LDPE has the greatest, LLDPE slightly less, and HDPE the least transparency. Transparency is reduced by crystallites if they are larger than the wavelength of visible light.
Manufacturing process
Monomer
The ingredient or monomer is ethylene (IUPAC name ethene), a gaseous hydrocarbon with the formula C2H4, which can be viewed as a pair of methylene groups (−−) connected to each other. Typical specifications for PE purity are <5 ppm for water, oxygen, and other alkenes contents. Acceptable contaminants include N2, ethane (common precursor to ethylene), and methane. Ethylene is usually produced from petrochemical sources, but is also generated by dehydration of ethanol.
Polymerization
Polymerization of ethylene to polyethylene is described by the following chemical equation:
(gas) → − (solid)Δ/ =
Ethylene is a stable molecule that polymerizes only upon contact with catalysts. The conversion is highly exothermic. Coordination polymerization is the most pervasive technology, which means that metal chlorides or metal oxides are used. The most common catalysts consist of titanium(III) chloride, the so-called Ziegler–Natta catalysts. Another common catalyst is the Phillips catalyst, prepared by depositing chromium(VI) oxide on silica. Polyethylene can be produced through radical polymerization, but this route has only limited utility and typically requires high-pressure apparatus.
Joining
Commonly used methods for joining polyethylene parts together include:
Welding
Hot gas welding
Infrared welding
Laser welding
Ultrasonic welding
Heat sealing
Heat fusion
Fastening
Adhesives
Pressure-sensitive adhesive (PSAs)
Dispersion of solvent-type PSAs
Polyurethane contact adhesives
Two-part polyurethane
Epoxy adhesives
Hot-melt adhesives
Solvent bonding – Adhesives and solvents are rarely used as solvent bonding because polyethylene is nonpolar and has a high resistance to solvents.
Pressure-sensitive adhesives (PSA) are feasible if the surface chemistry or charge is modified with plasma activation, flame treatment, or corona treatment.
Classification
Polyethylene is classified by its density and branching. Its mechanical properties depend significantly on variables such as the extent and type of branching, the crystal structure, and the molecular weight. There are several types of polyethylene:
Ultra-high-molecular-weight polyethylene (UHMWPE)
Ultra-low-molecular-weight polyethylene (ULMWPE or PE-WAX)
High-molecular-weight polyethylene (HMWPE)
High-density polyethylene (HDPE)
High-density cross-linked polyethylene (HDXLPE)
Cross-linked polyethylene (PEX or XLPE)
Medium-density polyethylene (MDPE)
Linear low-density polyethylene (LLDPE)
Low-density polyethylene (LDPE)
Very-low-density polyethylene (VLDPE)
Chlorinated polyethylene (CPE)
With regard to sold volumes, the most important polyethylene grades are HDPE, LLDPE, and LDPE.
Ultra-high-molecular-weight (UHMWPE)
UHMWPE is polyethylene with a molecular weight numbering in the millions, usually between 3.5 and 7.5 million amu. The high molecular weight makes it a very tough material, but results in less efficient packing of the chains into the crystal structure as evidenced by densities of less than high-density polyethylene (for example, 0.930–0.935 g/cm3). UHMWPE can be made through any catalyst technology, although Ziegler catalysts are most common. Because of its outstanding toughness and its cut, wear, and excellent chemical resistance, UHMWPE is used in a diverse range of applications. These include can- and bottle-handling machine parts, moving parts on weaving machines, bearings, gears, artificial joints, edge protection on ice rinks, steel cable replacements on ships, and butchers' chopping boards. It is commonly used for the construction of articular portions of implants used for hip and knee replacements. As fiber, it competes with aramid in bulletproof vests.
High-density (HDPE)
HDPE is defined by a density of greater or equal to 0.941 g/cm3. HDPE has a low degree of branching. The mostly linear molecules pack together well, so intermolecular forces are stronger than in highly branched polymers. HDPE can be produced by chromium/silica catalysts, Ziegler–Natta catalysts or metallocene catalysts; by choosing catalysts and reaction conditions, the small amount of branching that does occur can be controlled. These catalysts prefer the formation of free radicals at the ends of the growing polyethylene molecules. They cause new ethylene monomers to add to the ends of the molecules, rather than along the middle, causing the growth of a linear chain.
HDPE has high tensile strength. It is used in products and packaging such as milk jugs, detergent bottles, butter tubs, garbage containers, and water pipes.
Cross-linked (PEX or XLPE)
PEX is a medium- to high-density polyethylene containing cross-link bonds introduced into the polymer structure, changing the thermoplastic into a thermoset. The high-temperature properties of the polymer are improved, its flow is reduced, and its chemical resistance is enhanced. PEX is used in some potable-water plumbing systems because tubes made of the material can be expanded to fit over a metal nipple and it will slowly return to its original shape, forming a permanent, water-tight connection.
Medium-density (MDPE)
MDPE is defined by a density range of 0.926–0.940 g/cm3. MDPE can be produced by chromium/silica catalysts, Ziegler–Natta catalysts, or metallocene catalysts. MDPE has good shock and drop resistance properties. It also is less notch-sensitive than HDPE; stress-cracking resistance is better than HDPE. MDPE is typically used in gas pipes and fittings, sacks, shrink film, packaging film, carrier bags, and screw closures.
Linear low-density (LLDPE)
LLDPE is defined by a density range of 0.915–0.925 g/cm3. LLDPE is a substantially linear polymer with significant numbers of short branches, commonly made by copolymerization of ethylene with short-chain alpha-olefins (for example, 1-butene, 1-hexene, and 1-octene). LLDPE has higher tensile strength than LDPE, and it exhibits higher impact and puncture resistance than LDPE. Lower-thickness (gauge) films can be blown, compared with LDPE, with better environmental stress cracking resistance, but they are not as easy to process. LLDPE is used in packaging, particularly film for bags and sheets. Lower thickness may be used compared to LDPE. It is used for cable coverings, toys, lids, buckets, containers, and pipe. While other applications are available, LLDPE is used predominantly in film applications due to its toughness, flexibility, and relative transparency. Product examples range from agricultural films, Saran wrap, and bubble wrap to multilayer and composite films.
Low-density (LDPE)
LDPE is defined by a density range of 0.910–0.940 g/cm3. LDPE has a high degree of short- and long-chain branching, which means that the chains do not pack into the crystal structure as well. It has, therefore, less strong intermolecular forces as the instantaneous-dipole induced-dipole attraction is less. This results in a lower tensile strength and increased ductility. LDPE is created by free-radical polymerization. The high degree of branching with long chains gives molten LDPE unique and desirable flow properties. LDPE is used for both rigid containers and plastic film applications such as plastic bags and film wrap.
The radical polymerization process used to make LDPE does not include a catalyst that "supervises" the radical sites on the growing PE chains. (In HDPE synthesis, the radical sites are at the ends of the PE chains, because the catalyst stabilizes their formation at the ends.) Secondary radicals (in the middle of a chain) are more stable than primary radicals (at the end of the chain), and tertiary radicals (at a branch point) are more stable yet. Each time an ethylene monomer is added, it creates a primary radical, but often these will rearrange to form more stable secondary or tertiary radicals. Addition of ethylene monomers to the secondary or tertiary sites creates branching.
Very-low-density (VLDPE)
VLDPE is defined by a density range of 0.880–0.915 g/cm3. VLDPE is a substantially linear polymer with high levels of short-chain branches, commonly made by copolymerization of ethylene with short-chain alpha-olefins (for example, 1-butene, 1-hexene and 1-octene). VLDPE is most commonly produced using metallocene catalysts due to the greater co-monomer incorporation exhibited by these catalysts. VLDPEs are used for hose and tubing, ice and frozen food bags, food packaging and stretch wrap as well as impact modifiers when blended with other polymers.
Much research activity has focused on the nature and distribution of long chain branches in polyethylene. In HDPE, a relatively small number of these branches, perhaps one in 100 or 1,000 branches per backbone carbon, can significantly affect the rheological properties of the polymer.
Copolymers
In addition to copolymerization with alpha-olefins, ethylene can be copolymerized with a wide range of other monomers and ionic composition that creates ionized free radicals. Common examples include vinyl acetate (the resulting product is ethylene-vinyl acetate copolymer, or EVA, widely used in athletic-shoe sole foams) and a variety of acrylates. Applications of acrylic copolymer include packaging and sporting goods, and superplasticizer, used in cement production.
Types of polyethylenes
The particular material properties of "polyethylene" depend on its molecular structure. Molecular weight and crystallinity are the most significant factors; crystallinity in turn depends on molecular weight and degree of branching. The less the polymer chains are branched, and the lower the molecular weight, the higher the crystallinity of polyethylene. Crystallinity ranges from 35% (PE-LD/PE-LLD) to 80% (PE-HD). Polyethylene has a density of 1.0 g/cm3 in crystalline regions and 0.86 g/cm3 in amorphous regions. An almost linear relationship exists between density and crystallinity.
The degree of branching of the different types of polyethylene can be schematically represented as follows:
The figure shows polyethylene backbones, short-chain branches and side-chain branches. The polymer chains are represented linearly.
Chain branches
The properties of polyethylene are highly dependent on type and number of chain branches. The chain branches in turn depend on the process used: either the high-pressure process (only PE-LD) or the low-pressure process (all other PE grades). Low-density polyethylene is produced by the high-pressure process by radical polymerization, thereby numerous short chain branches as well as long chain branches are formed. Short chain branches are formed by intramolecular chain transfer reactions, they are always butyl or ethyl chain branches because the reaction proceeds after the following mechanism:
Environmental issues
The widespread usage of polyethylene poses potential difficulties for waste management because it is not readily biodegradable. Since 2008, Japan has increased plastic recycling, but still has a large amount of plastic wrapping which goes to waste. Plastic recycling in Japan is a potential US$90 billion market.
It is possible to rapidly convert polyethylene to hydrogen and graphene by heating. The energy needed is much less than for producing hydrogen by electrolysis.
Biodegradability
Several experiments have been conducted aimed at discovering enzyme or organisms that will degrade polyethylene. Several plastics - polyesters, polycarbonates, polyamides - degrade either by hydrolysis or air oxidation. In some cases the degradation is increased by bacteria or various enzyme cocktails. The situation is very different with polymers where the backbone consists solely of C-C bonds. These polymers include polyethylene, but also polypropylene, polystyrene and acrylates. At best, these polymers degrade very slowly, but these experiments are difficult because yields and rates are very slow. Further confusing the situation, even preliminary successes are greeted with enthusiasm by the popular press. Some technical challenges in this area include the failure to identify enzymes responsible for the proposed degradation. Another issue is that organisms are incapable of importing hydrocarbons of molecular weight greater than 500.
Bacteria and insect case studies
The Indian mealmoth larvae are claimed to metabolize polyethylene based on observing that plastic bags at a researcher's home had small holes in them. Deducing that the hungry larvae must have digested the plastic somehow, he and his team analyzed their gut bacteria and found a few that could use plastic as their only carbon source. Not only could the bacteria from the guts of the Plodia interpunctella moth larvae metabolize polyethylene, they degraded it significantly, dropping its tensile strength by 50%, its mass by 10% and the molecular weights of its polymeric chains by 13%.
The caterpillar of Galleria mellonella is claimed to consume polyethylene. The caterpillar is able to digest polyethylene due to a combination of its gut microbiota and its saliva containing enzymes that oxidise and depolymerise the plastic.
Climate change
When exposed to ambient solar radiation the plastic produces trace amounts of two greenhouse gases, methane and ethylene. The plastic type which releases gases at the highest rate is low-density polyethylene (LDPE). Due to its low density it breaks down more easily over time, leading to higher surface areas. When incubated in air, LDPE emits gases at rates ~2 times and ~76 times higher in comparison to incubation in water for methane and ethylene, respectively. However, based on the rates measured in the study methane production by plastics is presently an insignificant component of the global methane budget.
Chemically modified polyethylene
Polyethylene may either be modified in the polymerization by polar or non-polar comonomers or after polymerization through polymer-analogous reactions. Common polymer-analogous reactions are in case of polyethylene crosslinking, chlorination and sulfochlorination.
Non-polar ethylene copolymers
α-olefins
In the low pressure process α-olefins (e.g. 1-butene or 1-hexene) may be added, which are incorporated in the polymer chain during polymerization. These copolymers introduce short side chains, thus crystallinity and density are reduced. As explained above, mechanical and thermal properties are changed thereby. In particular, PE-LLD is produced this way.
Metallocene polyethylene (PE-MC)
Metallocene polyethylene (PE-M) is prepared by means of metallocene catalysts, usually including copolymers (z. B. ethene / hexene). Metallocene polyethylene has a relatively narrow molecular weight distribution, exceptionally high toughness, excellent optical properties and a uniform comonomer content. Because of the narrow molecular weight distribution it behaves less pseudoplastic (especially under larger shear rates). Metallocene polyethylene has a low proportion of low molecular weight (extractable) components and a low welding and sealing temperature. Thus, it is particularly suitable for the food industry.
Polyethylene with multimodal molecular weight distribution
Polyethylene with multimodal molecular weight distribution consists of several polymer fractions, which are homogeneously mixed. Such polyethylene types offer extremely high stiffness, toughness, strength, stress crack resistance and an increased crack propagation resistance. They consist of equal proportions higher and lower molecular polymer fractions. The lower molecular weight units crystallize easier and relax faster. The higher molecular weight fractions form linking molecules between crystallites, thereby increasing toughness and stress crack resistance. Polyethylene with multimodal molecular weight distribution can be prepared either in two-stage reactors, by catalysts with two active centers on a carrier or by blending in extruders.
Cyclic olefin copolymers (COC)
Cyclic olefin copolymers are prepared by copolymerization of ethene and cycloolefins (usually norbornene) produced by using metallocene catalysts. The resulting polymers are amorphous polymers and particularly transparent and heat resistant.
Polar ethylene copolymers
The basic compounds used as polar comonomers are vinyl alcohol (Ethenol, an unsaturated alcohol), acrylic acid (propenoic acid, an unsaturated acid) and esters containing one of the two compounds.
Ethylene copolymers with unsaturated alcohols
Ethylene/vinyl alcohol copolymer (EVOH) is (formally) a copolymer of PE and vinyl alcohol (ethenol), which is prepared by (partial) hydrolysis of ethylene-vinyl acetate copolymer (as vinyl alcohol itself is not stable). However, typically EVOH has a higher comonomer content than the VAC commonly used.
EVOH is used in multilayer films for packaging as a barrier layer (barrier plastic). As EVOH is hygroscopic (water-attracting), it absorbs water from the environment, whereby it loses its barrier effect. Therefore, it must be used as a core layer surrounded by other plastics (like LDPE, PP, PA or PET). EVOH is also used as a coating agent against corrosion at street lights, traffic light poles and noise protection walls.
Ethylene/acrylic acid copolymers (EAA)
Copolymer of ethylene and unsaturated carboxylic acids (such as acrylic acid) are characterized by good adhesion to diverse materials, by resistance to stress cracking and high flexibility. However, they are more sensitive to heat and oxidation than ethylene homopolymers. Ethylene/acrylic acid copolymers are used as adhesion promoters.
If salts of an unsaturated carboxylic acid are present in the polymer, thermo-reversible ion networks are formed, they are called ionomers. Ionomers are highly transparent thermoplastics which are characterized by high adhesion to metals, high abrasion resistance and high water absorption.
Ethylene copolymers with unsaturated esters
If unsaturated esters are copolymerized with ethylene, either the alcohol moiety may be in the polymer backbone (as it is the case in ethylene-vinyl acetate copolymer) or of the acid moiety (e. g. in ethylene-ethyl acrylate copolymer). Ethylene-vinyl acetate copolymers are prepared similarly to LD-PE by high pressure polymerization. The proportion of comonomer has a decisive influence on the behaviour of the polymer.
The density decreases up to a comonomer share of 10% because of the disturbed crystal formation. With higher proportions it approaches to the one of polyvinyl acetate (1.17 g/cm3). Due to decreasing crystallinity ethylene vinyl acetate copolymers are getting softer with increasing comonomer content. The polar side groups change the chemical properties significantly (compared to polyethylene): weather resistance, adhesiveness and weldability rise with comonomer content, while the chemical resistance decreases. Also mechanical properties are changed: stress cracking resistance and toughness in the cold rise, whereas yield stress and heat resistance decrease. With a very high proportion of comonomers (about 50%) rubbery thermoplastics are produced (thermoplastic elastomers).
Ethylene-ethyl acrylate copolymers behave similarly to ethylene-vinyl acetate copolymers.
Crosslinking
A basic distinction is made between peroxide crosslinking (PE-Xa), silane crosslinking (PE-Xb), electron beam crosslinking (PE-Xc) and azo crosslinking (PE-Xd).
Shown are the peroxide, the silane and irradiation crosslinking. In each method, a radical is generated in the polyethylene chain (top center), either by radiation (h·ν) or by peroxides (R-O-O-R). Then, two radical chains can either directly crosslink (bottom left) or indirectly by silane compounds (bottom right).
Peroxide crosslinking (PE-Xa): The crosslinking of polyethylene using peroxides (e. g. dicumyl or di-tert-butyl peroxide) is still of major importance. In the so-called Engel process, a mixture of HDPE and 2% peroxide is at first mixed at low temperatures in an extruder and then crosslinked at high temperatures (between 200 and 250 °C). The peroxide decomposes to peroxide radicals (RO•), which abstract (remove) hydrogen atoms from the polymer chain, leading to radicals. When these combine, a crosslinked network is formed. The resulting polymer network is uniform, of low tension and high flexibility, whereby it is softer and tougher than (the irradiated) PE-Xc.
Silane crosslinking (PE-Xb): In the presence of silanes (e.g. trimethoxyvinylsilane) polyethylene can initially be Si-functionalized by irradiation or by a small amount of a peroxide. Later Si-OH groups can be formed in a water bath by hydrolysis, which condense then and crosslink the PE by the formation of Si-O-Si bridges. [16] Catalysts such as dibutyltin dilaurate may accelerate the reaction.
Irradiation crosslinking (PE-Xc): The crosslinking of polyethylene is also possible by a downstream radiation source (usually an electron accelerator, occasionally an isotopic radiator). PE products are crosslinked below the crystalline melting point by splitting off hydrogen atoms. β-radiation possesses a penetration depth of 10 mm, ɣ-radiation 100 mm. Thereby the interior or specific areas can be excluded from the crosslinking. However, due to high capital and operating costs radiation crosslinking plays only a minor role compared with the peroxide crosslinking. In contrast to peroxide crosslinking, the process is carried out in the solid state. Thereby, the cross-linking takes place primarily in the amorphous regions, while the crystallinity remains largely intact.
Azo crosslinking (PE-Xd): In the so-called Lubonyl process polyethylene is crosslinked preadded azo compounds after extrusion in a hot salt bath.
Chlorination and sulfochlorination
Chlorinated Polyethylene (PE-C) is an inexpensive material having a chlorine content from 34 to 44%. It is used in blends with PVC because the soft, rubbery chloropolyethylene is embedded in the PVC matrix, thereby increasing the impact resistance. It also increases the weather resistance. Furthermore, it is used for softening PVC foils, without risking the migrate of plasticizers. Chlorinated polyethylene can be crosslinked peroxidically to form an elastomer which is used in cable and rubber industry. When chlorinated polyethylene is added to other polyolefins, it reduces the flammability.
Chlorosulfonated PE (CSM) is used as starting material for ozone-resistant synthetic rubber.
Bio-based polyethylene
Braskem and Toyota Tsusho Corporation started joint marketing activities to produce polyethylene from sugarcane. Braskem will build a new facility at their existing industrial unit in Triunfo, Rio Grande do Sul, Brazil with an annual production capacity of , and will produce high-density and low-density polyethylene from bioethanol derived from sugarcane.
Nomenclature and general description of the process
The name polyethylene comes from the ingredient and not the resulting chemical compound, which contains no double bonds. The scientific name polyethene is systematically derived from the scientific name of the monomer. The alkene monomer converts to a long, sometimes very long, alkane in the polymerization process. In certain circumstances it is useful to use a structure-based nomenclature; in such cases IUPAC recommends poly(methylene) (poly(methanediyl) is a non-preferred alternative). The difference in names between the two systems is due to the opening up of the monomer's double bond upon polymerization. The name is abbreviated to PE. In a similar manner polypropylene and polystyrene are shortened to PP and PS, respectively. In the United Kingdom and India the polymer is commonly called polythene, from the ICI trade name, although this is not recognized scientifically.
| Physical sciences | Hydrocarbons | null |
77432 | https://en.wikipedia.org/wiki/Hypertension | Hypertension | Hypertension, also known as high blood pressure, is a long-term medical condition in which the blood pressure in the arteries is persistently elevated. High blood pressure usually does not cause symptoms itself. It is, however, a major risk factor for stroke, coronary artery disease, heart failure, atrial fibrillation, peripheral arterial disease, vision loss, chronic kidney disease, and dementia. Hypertension is a major cause of premature death worldwide.
High blood pressure is classified as primary (essential) hypertension or secondary hypertension. About 90–95% of cases are primary, defined as high blood pressure due to nonspecific lifestyle and genetic factors. Lifestyle factors that increase the risk include excess salt in the diet, excess body weight, smoking, physical inactivity and alcohol use. The remaining 5–10% of cases are categorized as secondary hypertension, defined as high blood pressure due to a clearly identifiable cause, such as chronic kidney disease, narrowing of the kidney arteries, an endocrine disorder, or the use of birth control pills.
Blood pressure is classified by two measurements, the systolic (first number) and diastolic (second number) pressures. For most adults, normal blood pressure at rest is within the range of 100–140 millimeters mercury (mmHg) systolic and 60–90 mmHg diastolic. For most adults, high blood pressure is present if the resting blood pressure is persistently at or above 130/80 or 140/90 mmHg. Different numbers apply to children. Ambulatory blood pressure monitoring over a 24-hour period appears more accurate than office-based blood pressure measurement.
Lifestyle changes and medications can lower blood pressure and decrease the risk of health complications. Lifestyle changes include weight loss, physical exercise, decreased salt intake, reducing alcohol intake, and a healthy diet. If lifestyle changes are not sufficient, blood pressure medications are used. Up to three medications taken concurrently can control blood pressure in 90% of people. The treatment of moderately high arterial blood pressure (defined as >160/100 mmHg) with medications is associated with an improved life expectancy. The effect of treatment of blood pressure between 130/80 mmHg and 160/100 mmHg is less clear, with some reviews finding benefit and others finding unclear benefit. High blood pressure affects 33% of the population globally. About half of all people with high blood pressure do not know that they have it. In 2019, high blood pressure was believed to have been a factor in 19% of all deaths (10.4 million globally).
Signs and symptoms
Hypertension is rarely accompanied by symptoms. Half of all people with hypertension are unaware that they have it. Hypertension is usually identified as part of health screening or when seeking healthcare for an unrelated problem.
Some people with high blood pressure report headaches, as well as lightheadedness, vertigo, tinnitus (buzzing or hissing in the ears), altered vision or fainting episodes. These symptoms, however, might be related to associated anxiety rather than the high blood pressure itself.
Long-standing untreated hypertension can cause organ damage with signs such as changes in the optic fundus seen by ophthalmoscopy. The severity of hypertensive retinopathy correlates roughly with the duration or the severity of the hypertension. Other hypertension-caused organ damage include chronic kidney disease and thickening of the heart muscle.
Secondary hypertension
Secondary hypertension is hypertension due to an identifiable cause, and may result in certain specific additional signs and symptoms. For example, as well as causing high blood pressure, Cushing's syndrome frequently causes truncal obesity, glucose intolerance, moon face, a hump of fat behind the neck and shoulders (referred to as a buffalo hump), and purple abdominal stretch marks. Hyperthyroidism frequently causes weight loss with increased appetite, fast heart rate, bulging eyes, and tremor. Renal artery stenosis may be associated with a localized abdominal bruit to the left or right of the midline, or in both locations. Coarctation of the aorta frequently causes a decreased blood pressure in the lower extremities relative to the arms, or delayed or absent femoral arterial pulses. Pheochromocytoma may cause abrupt episodes of hypertension accompanied by headache, palpitations, pale appearance, and excessive sweating.
Hypertensive crisis
Severely elevated blood pressure (equal to or greater than a systolic 180 mmHg or diastolic of 120 mmHg) is referred to as a hypertensive crisis. Hypertensive crisis is categorized as either hypertensive urgency or hypertensive emergency, according to the absence or presence of end organ damage, respectively.
In hypertensive urgency, there is no evidence of end organ damage resulting from the elevated blood pressure. In these cases, oral medications are used to lower the BP gradually over 24 to 48 hours.
In hypertensive emergency, there is evidence of direct damage to one or more organs. The most affected organs include the brain, kidney, heart and lungs, producing symptoms which may include confusion, drowsiness, chest pain and breathlessness. In hypertensive emergency, the blood pressure must be reduced more rapidly to stop ongoing organ damage; however, there is a lack of randomized controlled trial evidence for this approach.
Pregnancy
Hypertension occurs in approximately 8–10% of pregnancies. Two blood pressure measurements six hours apart of greater than 140/90 mmHg are diagnostic of hypertension in pregnancy. High blood pressure in pregnancy can be classified as pre-existing hypertension, gestational hypertension, or pre-eclampsia. Women who have chronic hypertension before their pregnancy are at increased risk of complications such as premature birth, low birthweight or stillbirth. Women who have high blood pressure and had complications in their pregnancy have three times the risk of developing cardiovascular disease compared to women with normal blood pressure who had no complications in pregnancy.
Pre-eclampsia is a serious condition of the second half of pregnancy and following delivery characterised by increased blood pressure and the presence of protein in the urine. It occurs in about 5% of pregnancies and is responsible for approximately 16% of all maternal deaths globally. Pre-eclampsia also doubles the risk of death of the baby around the time of birth. Usually there are no symptoms in pre-eclampsia and it is detected by routine screening. When symptoms of pre-eclampsia occur the most common are headache, visual disturbance (often "flashing lights"), vomiting, pain over the stomach, and swelling. Pre-eclampsia can occasionally progress to a life-threatening condition called eclampsia, which is a hypertensive emergency and has several serious complications including vision loss, brain swelling, seizures, kidney failure, pulmonary edema, and disseminated intravascular coagulation (a blood clotting disorder).
In contrast, gestational hypertension is defined as new-onset hypertension during pregnancy without protein in the urine.
There have been significant findings on how exercising can help reduce the effects of hypertension just after one bout of exercise. Exercising can help reduce hypertension as well as pre-eclampsia and eclampsia.
The acute physiological responses include an increase in cardiac output (CO) of the individual (increased heart rate and stroke volume). This increase in CO can inadvertently maintain the amount of blood going into the muscles, improving functionality of the muscle later. Exercising can also improve systolic and diastolic blood pressure making it easier for blood to pump to the body. Through regular bouts of physical activity, blood pressure can reduce the incidence of hypertension.
Aerobic exercise has been shown to regulate blood pressure more effectively than resistance training. It is recommended to see the effects of exercising, that a person should aim for 5-7 days/ week of aerobic exercise. This type of exercise should have an intensity of light to moderate, utilizing ~85% of max heart rate (220-age). Aerobic has shown a decrease in SBP by 5-15mmHg, versus resistance training showing a decrease of only 3-5mmHg. Aerobic exercises such as jogging, rowing, dancing, or hiking can decrease SBP the greatest. The decrease in SBP can regulate the effect of hypertension ensuring the baby will not be harmed. Resistance training takes a toll on the cardiovascular system in untrained individuals, leading to a reluctance in prescription of resistance training for hypertensive reduction purposes.
Children
Failure to thrive, seizures, irritability, lack of energy, and difficulty in breathing can be associated with hypertension in newborns and young infants. In older infants and children, hypertension can cause headache, unexplained irritability, fatigue, failure to thrive, blurred vision, nosebleeds, and facial paralysis.
Causes
Primary hypertension
Primary (also termed essential) hypertension results from a complex interaction of genes and environmental factors. More than 2000 common genetic variants with small effects on blood pressure have been identified in association with high blood pressure, as well as some rare genetic variants with large effects on blood pressure. There is also evidence that DNA methylation at multiple nearby CpG sites may link some sequence variation to blood pressure, possibly via effects on vascular or renal function.
Blood pressure rises with aging in societies with a western diet and lifestyle, and the risk of becoming hypertensive in later life is substantial in most such societies. Several environmental or lifestyle factors influence blood pressure. Reducing dietary salt intake lowers blood pressure; as does weight loss, exercise training, vegetarian diets, increased dietary potassium intake and high dietary calcium supplementation. Increasing alcohol intake is associated with higher blood pressure, but the possible roles of other factors such as caffeine consumption, and vitamin D deficiency are less clear. Average blood pressure is higher in the winter than in the summer.
Depression is associated with hypertension and loneliness is also a risk factor. Periodontal disease is also associated with high blood pressure. Arsenic exposure through drinking water is associated with elevated blood pressure. Air pollution is associated with hypertension. Whether these associations are causal is unknown. Gout and elevated blood uric acid are associated with hypertension and evidence from genetic (Mendelian Randomization) studies and clinical trials indicate this relationship is likely to be causal. Insulin resistance, which is common in obesity and is a component of syndrome X (or metabolic syndrome), can cause hyperuricemia and gout and is also associated with elevated blood pressure.
Events in early life, such as low birth weight, maternal smoking, and lack of breastfeeding may be risk factors for adult essential hypertension, although strength of the relationships is weak and the mechanisms linking these exposures to adult hypertension remain unclear.
Secondary hypertension
Secondary hypertension results from an identifiable cause. Kidney disease is the most common secondary cause of hypertension. Hypertension can also be caused by endocrine conditions, such as Cushing's syndrome, hyperthyroidism, hypothyroidism, acromegaly, Conn's syndrome or hyperaldosteronism, renal artery stenosis (from atherosclerosis or fibromuscular dysplasia), hyperparathyroidism, and pheochromocytoma. Other causes of secondary hypertension include obesity, sleep apnea, pregnancy, coarctation of the aorta, excessive eating of liquorice, excessive drinking of alcohol, certain prescription medicines, herbal remedies, and stimulants such as cocaine and methamphetamine.
A 2018 review found that any alcohol increased blood pressure in males while over one or two drinks increased the risk in females.
Pathophysiology
In most people with established essential hypertension, increased resistance to blood flow (total peripheral resistance) accounts for the high pressure while cardiac output remains normal. There is evidence that some younger people with prehypertension or 'borderline hypertension' have high cardiac output, an elevated heart rate and normal peripheral resistance, termed hyperkinetic borderline hypertension. These individuals may develop the typical features of established essential hypertension in later life as their cardiac output falls and peripheral resistance rises with age. Whether this pattern is typical of all people who ultimately develop hypertension is disputed. The increased peripheral resistance in established hypertension is mainly attributable to structural narrowing of small arteries and arterioles, although a reduction in the number or density of capillaries may also contribute.
It is not clear whether or not vasoconstriction of arteriolar blood vessels plays a role in hypertension. Hypertension is also associated with decreased peripheral venous compliance, which may increase venous return, increase cardiac preload and, ultimately, cause diastolic dysfunction. For patients having hypertension, higher heart rate variability (HRV) is a risk factor for atrial fibrillation.
Pulse pressure (the difference between systolic and diastolic blood pressure) is frequently increased in older people with hypertension. This can mean that systolic pressure is abnormally high, but diastolic pressure may be normal or low, a condition termed isolated systolic hypertension. The high pulse pressure in elderly people with hypertension or isolated systolic hypertension is explained by increased arterial stiffness, which typically accompanies aging and may be exacerbated by high blood pressure.
Many mechanisms have been proposed to account for the rise in peripheral resistance in hypertension. Most evidence implicates either disturbances in the kidneys' salt and water handling (particularly abnormalities in the intrarenal renin–angiotensin system) or abnormalities of the sympathetic nervous system. These mechanisms are not mutually exclusive and it is likely that both contribute to some extent in most cases of essential hypertension. It has also been suggested that endothelial dysfunction and vascular inflammation may also contribute to increased peripheral resistance and vascular damage in hypertension. Interleukin 17 has garnered interest for its role in increasing the production of several other immune system chemical signals thought to be involved in hypertension such as tumor necrosis factor alpha, interleukin 1, interleukin 6, and interleukin 8.
Excessive sodium or insufficient potassium in the diet leads to excessive intracellular sodium, which contracts vascular smooth muscle, restricting blood flow and so increases blood pressure. Non-modulating essential hypertension is a form of salt-sensitive hypertension, where sodium intake does not modulate either adrenal or renal vascular responses to angiotensin II. They make up 25% of the hypertensive population.
Diagnosis
Hypertension is diagnosed on the basis of a persistently high resting blood pressure. Elevated blood pressure measurements on at least two separate occasions is required for a diagnosis of hypertension.
Measurement technique
For an accurate diagnosis of hypertension to be made, it is essential for proper blood pressure measurement technique to be used. Improper measurement of blood pressure is common and can change the blood pressure reading by up to 10 mmHg, which can lead to misdiagnosis and misclassification of hypertension. Correct blood pressure measurement technique involves several steps. Proper blood pressure measurement requires the person whose blood pressure is being measured to sit quietly for at least five minutes which is then followed by application of a properly fitted blood pressure cuff to a bare upper arm. The person should be seated with their back supported, feet flat on the floor, and with their legs uncrossed. The person whose blood pressure is being measured should avoid talking or moving during this process. The arm being measured should be supported on a flat surface at the level of the heart. Blood pressure measurement should be done in a quiet room so the medical professional checking the blood pressure can hear the Korotkoff sounds while listening to the brachial artery with a stethoscope for accurate blood pressure measurements. The blood pressure cuff should be deflated slowly (2–3 mmHg per second) while listening for the Korotkoff sounds. The bladder should be emptied before a person's blood pressure is measured since this can increase blood pressure by up to 15/10 mmHg. Multiple blood pressure readings (at least two) spaced 1–2 minutes apart should be obtained to ensure accuracy. Ambulatory blood pressure monitoring over 12 to 24 hours is the most accurate method to confirm the diagnosis. An exception to this is those with very high blood pressure readings especially when there is poor organ function.
With the availability of 24-hour ambulatory blood pressure monitors and home blood pressure machines, the importance of not wrongly diagnosing those who have white coat hypertension has led to a change in protocols. In the United Kingdom, current best practice is to follow up a single raised clinic reading with ambulatory measurement, or less ideally with home blood pressure monitoring over the course of 7 days. The United States Preventive Services Task Force also recommends getting measurements outside of the healthcare environment. Pseudohypertension in the elderly or noncompressibility artery syndrome may also require consideration. This condition is believed to be due to calcification of the arteries resulting in abnormally high blood pressure readings with a blood pressure cuff while intra arterial measurements of blood pressure are normal. Orthostatic hypertension is when blood pressure increases upon standing.
Other investigations
Once the diagnosis of hypertension has been made, further testing may be performed to find secondary hypertension, identify comorbidities such as diabetes, identify hypertension-caused organ damage such as chronic kidney disease or thickening of the heart muscle, and for cardiovascular disease risk stratification.
Secondary hypertension is more common in preadolescent children, with most cases caused by kidney disease. Primary or essential hypertension is more common in adolescents and adults and has multiple risk factors, including obesity and a family history of hypertension.
Initial assessment upon diagnosis of hypertension should include a complete history and physical examination. The World Health Organization suggests the following initial tests: serum electrolytes, serum creatinine, lipid panel, HbA1c or fasting glucose, urine dipstick and electrocardiogram (ECG/EKG). Serum creatinine is measured to assess for the presence of kidney disease, which can be either the cause or the result of hypertension. eGFR can also provide a baseline measurement of kidney function that can be used to monitor for side effects of certain anti-hypertensive drugs on kidney function. Testing of urine samples for protein is used as a secondary indicator of kidney disease. Lipid panel and glucose tests are done to identify comorbidities such as diabetes and hyperlipidemia and for cardiovascular risk stratification. Electrocardiogram (EKG/ECG) testing is done to check for evidence that the heart is under strain from high blood pressure, such as thickening of the heart muscle or whether the heart has experienced a prior minor disturbance such as a silent heart attack.
Classification in adults
Blood pressure measurements can be influenced by circumstances of measurement. Guidelines use different thresholds for office (also known as clinic), home (when the patient measures their own blood pressure at home), and ambulatory blood pressure (using an automated device over a 24-hour period).
Children
Hypertension occurs in around 0.2 to 3% of newborns; however, blood pressure is not measured routinely in healthy newborns. Hypertension is more common in high risk newborns. A variety of factors, such as gestational age, postconceptional age and birth weight needs to be taken into account when deciding if a blood pressure is normal in a newborn.
Hypertension defined as elevated blood pressure over several visits affects 1% to 5% of children and adolescents and is associated with long-term risks of ill-health. Blood pressure rises with age in childhood and, in children, hypertension is defined as an average systolic or diastolic blood pressure on three or more occasions equal or higher than the 95th percentile appropriate for the sex, age and height of the child. High blood pressure must be confirmed on repeated visits however before characterizing a child as having hypertension. In adolescents, it has been proposed that hypertension is diagnosed and classified using the same criteria as in adults.
Prevention
Much of the disease burden of high blood pressure is experienced by people who are not labeled as hypertensive. Consequently, population strategies are required to reduce the consequences of high blood pressure and reduce the need for antihypertensive medications. Lifestyle changes are recommended to lower blood pressure.
Recommended lifestyle changes for the prevention of hypertension include:
maintain normal body weight for adults (e.g. body mass index below 25 kg/m2)
reduce dietary sodium intake to <100 mmol/day (<6 g of salt (sodium chloride) or <2.4 g of sodium per day)
engage in regular aerobic physical activity with moderate intensity (minimum 150 minutes per week)
limit alcohol consumption, max 1 drink for women and 2 for men per day
consume a diet rich in whole grains, fruit and vegetables, such as the DASH diet
not smoking
stress reduction and management, e.g. by meditation and yoga
Effective lifestyle modification may lower blood pressure as much as an individual antihypertensive medication. Combinations of two or more lifestyle modifications can achieve even better results. There is considerable evidence that reducing dietary salt intake lowers blood pressure, but whether this translates into a reduction in mortality and cardiovascular disease remains uncertain. Estimated sodium intake ≥6 g/day and <3 g/day are both associated with high risk of death or major cardiovascular disease, but the association between high sodium intake and adverse outcomes is only observed in people with hypertension. Consequently, in the absence of results from randomized controlled trials, the wisdom of reducing levels of dietary sodium intake below 3 g/day has been questioned. ESC guidelines mention periodontitis is associated with poor cardiovascular health status.
The value of routine screening for hypertension is debated. In 2004, the National High Blood Pressure Education Program recommended that children aged 3 years and older have blood pressure measurement at least once at every health care visit and the National Heart, Lung, and Blood Institute and American Academy of Pediatrics made a similar recommendation. However, the American Academy of Family Physicians supports the view of the U.S. Preventive Services Task Force that the available evidence is insufficient to determine the balance of benefits and harms of screening for hypertension in children and adolescents who do not have symptoms. The US Preventive Services Task Force recommends screening adults 18 years or older for hypertension with office blood pressure measurement.
Management
According to one review published in 2003, reduction of the blood pressure by 5 mmHg can decrease the risk of stroke by 34%, of ischemic heart disease by 21%, and reduce the likelihood of dementia, heart failure, and mortality from cardiovascular disease.
Target blood pressure
Various expert groups have produced guidelines regarding how low the blood pressure target should be when a person is treated for hypertension. These groups recommend a target below the range 140–160 / 90–100 mmHg for the general population. Cochrane reviews recommend similar targets for subgroups such as people with diabetes and people with prior cardiovascular disease. Additionally, Cochrane reviews have found that for older individuals with moderate to high cardiovascular risk, the benefits of trying to achieve a lower than standard blood pressure target (at or below 140/90 mmHg) are outweighed by the risk associated with the intervention. These findings may not be applicable to other populations.
Many expert groups recommend a slightly higher target of 150/90 mmHg for those over somewhere between 60 and 80 years of age. The JNC 8 and American College of Physicians recommend the target of 150/90 mmHg for those over 60 years of age, but some experts within these groups disagree with this recommendation. Some expert groups have also recommended slightly lower targets in those with diabetes or chronic kidney disease, but others recommend the same target as for the general population. The issue of what is the best target and whether targets should differ for high risk individuals is unresolved, although some experts propose more intensive blood pressure lowering than advocated in some guidelines.
For people who have never experienced cardiovascular disease who are at a 10-year risk of cardiovascular disease of less than 10%, the 2017 American Heart Association guidelines recommend medications if the systolic blood pressure is >140 mmHg or if the diastolic BP is >90 mmHg. For people who have experienced cardiovascular disease or those who are at a 10-year risk of cardiovascular disease of greater than 10%, it recommends medications if the systolic blood pressure is >130 mmHg or if the diastolic BP is >80 mmHg.
Lifestyle modifications
The first line of treatment for hypertension is lifestyle changes, including dietary changes, physical activity, and weight loss. Though these have all been recommended in scientific advisories, a Cochrane systematic review found no evidence (due to lack of data) for effects of weight loss diets on death, long-term complications or adverse events in persons with hypertension. The review did find a decrease in body weight and blood pressure. Their potential effectiveness is similar to and at times exceeds a single medication. If hypertension is high enough to justify immediate use of medications, lifestyle changes are still recommended in conjunction with medication.
Dietary changes shown to reduce blood pressure include diets with low sodium, the DASH diet (Dietary Approaches to Stop Hypertension), which was the best against 11 other diet in an umbrella review, and plant-based diets. A 2024 clinical guideline recommended an increase dietary fiber intake, with a minimum of 28g/day for women and 38g/day for men diagnosed with hypertension.
Increasing dietary potassium has a potential benefit for lowering the risk of hypertension. The 2015 Dietary Guidelines Advisory Committee (DGAC) stated that potassium is one of the shortfall nutrients which is under-consumed in the United States. However, people who take certain antihypertensive medications (such as ACE-inhibitors or ARBs) should not take potassium supplements or potassium-enriched salts due to the risk of high levels of potassium.
Physical exercise regimens which are shown to reduce blood pressure include isometric resistance exercise, aerobic exercise, resistance exercise, and device-guided breathing.
Stress reduction techniques such as biofeedback or transcendental meditation may be considered as an add-on to other treatments to reduce hypertension, but do not have evidence for preventing cardiovascular disease on their own. Self-monitoring and appointment reminders might support the use of other strategies to improve blood pressure control, but need further evaluation.
Medications
Several classes of medications, collectively referred to as antihypertensive medications, are available for treating hypertension.
First-line medications for hypertension include thiazide-diuretics, calcium channel blockers, angiotensin converting enzyme inhibitors (ACE inhibitors), and angiotensin receptor blockers (ARBs). These medications may be used alone or in combination (ACE inhibitors and ARBs are not recommended for use together); the latter option may serve to minimize counter-regulatory mechanisms that act to restore blood pressure values to pre-treatment levels, although the evidence for first-line combination therapy is not strong enough. Most people require more than one medication to control their hypertension. Medications for blood pressure control should be implemented by a stepped care approach when target levels are not reached. Withdrawal of such medications in the elderly can be considered by healthcare professionals, because there is no strong evidence of an effect on mortality, myocardial infarction, or stroke.
Previously, beta-blockers such as atenolol were thought to have similar beneficial effects when used as first-line therapy for hypertension. However, a Cochrane review that included 13 trials found that the effects of beta-blockers are inferior to that of other antihypertensive medications in preventing cardiovascular disease.
The prescription of antihypertensive medication for children with hypertension has limited evidence. There is limited evidence which compare it with placebo and shows modest effect to blood pressure in short term. Administration of higher dose did not make the reduction of blood pressure greater.
Resistant hypertension
Resistant hypertension is defined as high blood pressure that remains above a target level, in spite of being prescribed three or more antihypertensive drugs simultaneously with different mechanisms of action. Failing to take prescribed medications as directed is an important cause of resistant hypertension.
Some common secondary causes of resistant hypertension include obstructive sleep apnea, primary aldosteronism and renal artery stenosis, and some rare secondary causes are pheochromocytoma and coarctation of the aorta. As many as one in five people with resistant hypertension have primary aldosteronism, which is a treatable and sometimes curable condition. Resistant hypertension may also result from chronically high activity of the autonomic nervous system, an effect known as neurogenic hypertension. Electrical therapies that stimulate the baroreflex are being studied as an option for lowering blood pressure in people in this situation.
Refractory hypertension is described by one source as elevated blood pressure unmitigated by five or more concurrent antihypertensive agents of different classes. People with refractory hypertension typically have increased sympathetic nervous system activity, and are at high risk for more severe cardiovascular diseases and all-cause mortality.
Epidemiology
Adults
, one in three or 33% of the world population were estimated to have hypertension. Of all people with hypertension, about 46% do not have a diagnosis of hypertension and are unaware that they have the condition. In 1975, almost 600 million people had a diagnosis of hypertension, a number which increased to 1.13 billion by 2015 mostly due to risk factors for hypertension increasing in low- and middle-income countries.
Hypertension is slightly more frequent in men. In people aged under 50 years, more men than women have hypertension, and in ages above 50 years the prevalence of hypertension is the same in men and women. In ages above 65 years, more women than men have hypertension. Hypertension becomes more common with age. Hypertension is common in high, medium, and low-income countries. It is more common in people of low socioeconomic status. Hypertension is around twice as common in diabetics.
In 2019, rates of diagnosed hypertension were highest in Africa (30% for both sexes), and lowest in the Americas (18% for both sexes). Rates also vary markedly within regions with country-level rates as low as 22.8% (men) and 18.4% (women) in Peru and as high as 61.6% (men) and 50.9% (women) in Paraguay.
In 1995 it was estimated that 24% of the United States population had hypertension or were taking antihypertensive medication. By 2004 this had increased to 29% and further to 32% (76 million US adults) by 2017. In 2017, with the American guidelines' change in definition for hypertension, 46% of people in the United States are affected. Some data shows African-American adults in the United States have among the highest rates of hypertension in the world at 44%. However, other research argues there has been a "myopic perspective" on American data and notes that other groups, particularly Russians and Eastern Europeans, have markedly higher rates of hypertension than Black Americans. Differences in hypertension rates are multifactorial and under study.
Children
Rates of high blood pressure in children and adolescents have increased in the last 20 years in the United States. Childhood hypertension, particularly in pre-adolescents, is more often secondary to an underlying disorder than in adults. Kidney disease is the most common secondary cause of hypertension in children and adolescents. Nevertheless, primary or essential hypertension accounts for most cases.
Prognosis
Hypertension is the most important preventable risk factor for premature death worldwide. It increases the risk of ischemic heart disease, stroke, peripheral vascular disease, and other cardiovascular diseases, including heart failure, aortic aneurysms, diffuse atherosclerosis, chronic kidney disease, atrial fibrillation, cancers, leukemia and pulmonary embolism. Hypertension is also a risk factor for cognitive impairment and dementia. Other complications include hypertensive retinopathy and hypertensive nephropathy.
History
Measurement
Modern understanding of the cardiovascular system began with the work of physician William Harvey (1578–1657), who described the circulation of blood in his book "De motu cordis". The English clergyman Stephen Hales made the first published measurement of blood pressure in 1733. However, hypertension as a clinical entity came into its own with the invention of the cuff-based sphygmomanometer by Scipione Riva-Rocci in 1896. This allowed easy measurement of systolic pressure in the clinic. In 1905, Nikolai Korotkoff improved the technique by describing the Korotkoff sounds that are heard when the artery is ausculted with a stethoscope while the sphygmomanometer cuff is deflated. This permitted systolic and diastolic pressure to be measured.
Identification
The symptoms similar to symptoms of patients with hypertensive crisis are discussed in medieval Persian medical texts in the chapter of "fullness disease". The symptoms include headache, heaviness in the head, sluggish movements, general redness and warm to touch feel of the body, prominent, distended and tense vessels, fullness of the pulse, distension of the skin, coloured and dense urine, loss of appetite, weak eyesight, impairment of thinking, yawning, drowsiness, vascular rupture, and hemorrhagic stroke. Fullness disease was presumed to be due to an excessive amount of blood within the blood vessels.
Descriptions of hypertension as a disease came among others from Thomas Young in 1808 and especially Richard Bright in 1836. The first report of elevated blood pressure in a person without evidence of kidney disease was made by Frederick Akbar Mahomed (1849–1884).
Until the 1990s, systolic hypertension was defined as systolic blood pressure of 160 mm Hg or greater. In 1993, the WHO/ISH guidelines defined 140 mmHg as the threshold for hypertension.
Treatment
Historically the treatment for what was called the "hard pulse disease" consisted in reducing the quantity of blood by bloodletting or the application of leeches. This was advocated by The Yellow Emperor of China, Cornelius Celsus, Galen, and Hippocrates. The therapeutic approach for the treatment of hard pulse disease included changes in lifestyle (staying away from anger and sexual intercourse) and dietary program for patients (avoiding the consumption of wine, meat, and pastries, reducing the volume of food in a meal, maintaining a low-energy diet and the dietary usage of spinach and vinegar).
In the 19th and 20th centuries, before effective pharmacological treatment for hypertension became possible, three treatment modalities were used, all with numerous side-effects: strict sodium restriction (for example the rice diet), sympathectomy (surgical ablation of parts of the sympathetic nervous system), and pyrogen therapy (injection of substances that caused a fever, indirectly reducing blood pressure).
The first chemical for hypertension, sodium thiocyanate, was used in 1900 but had many side effects and was unpopular. Several other agents were developed after the Second World War, the most popular and reasonably effective of which were tetramethylammonium chloride, hexamethonium, hydralazine, and reserpine (derived from the medicinal plant Rauvolfia serpentina). None of these were well tolerated. A major breakthrough was achieved with the discovery of the first well-tolerated orally available agents. The first was chlorothiazide, the first thiazide diuretic and developed from the antibiotic sulfanilamide, which became available in 1958. Subsequently, beta blockers, calcium channel blockers, angiotensin converting enzyme (ACE) inhibitors, angiotensin receptor blockers, and renin inhibitors were developed as antihypertensive agents.
Society and culture
Awareness
The World Health Organization has identified hypertension (high blood pressure) as the leading cause of cardiovascular mortality. The World Hypertension League (WHL), an umbrella organization of 85 national hypertension societies and leagues, recognized that more than 50% of the hypertensive population worldwide are unaware of their condition. To address this problem, the WHL initiated a global awareness campaign on hypertension in 2005 and dedicated 17 May of each year as World Hypertension Day.
Economics
High blood pressure is the most common chronic medical problem prompting visits to primary health care providers in US. The American Heart Association estimated the direct and indirect costs of high blood pressure in 2010 as $76.6 billion. In the US 80% of people with hypertension are aware of their condition, 71% take some antihypertensive medication, but only 48% of people aware that they have hypertension adequately control it. Adequate management of hypertension can be hampered by inadequacies in the diagnosis, treatment, or control of high blood pressure. Health care providers face many obstacles to achieving blood pressure control, including resistance to taking multiple medications to reach blood pressure goals. People also face the challenges of adhering to medicine schedules and making lifestyle changes. Nonetheless, the achievement of blood pressure goals is possible, and most importantly, lowering blood pressure significantly reduces the risk of death due to heart disease and stroke, the development of other debilitating conditions, and the cost associated with advanced medical care.
Other animals
Hypertension in cats is indicated with a systolic blood pressure greater than 150 mmHg, with amlodipine the usual first-line treatment. A cat with a systolic blood pressure above 170 mmHg is considered hypertensive. If a cat has other problems such as any kidney disease or retina detachment then a blood pressure below 160 mmHg may also need to be monitored.
Normal blood pressure in dogs can differ substantially between breeds but hypertension is often diagnosed if systolic blood pressure is above 160 mmHg particularly if this is associated with target organ damage. Inhibitors of the renin-angiotensin system and calcium channel blockers are often used to treat hypertension in dogs, although other drugs may be indicated for specific conditions causing high blood pressure.
| Biology and health sciences | Non-infectious disease | null |
77473 | https://en.wikipedia.org/wiki/Flerovium | Flerovium | Flerovium is a synthetic chemical element; it has symbol Fl and atomic number 114. It is an extremely radioactive, superheavy element, named after the Flerov Laboratory of Nuclear Reactions of the Joint Institute for Nuclear Research in Dubna, Russia, where the element was discovered in 1999. The lab's name, in turn, honours Russian physicist Georgy Flyorov ( in Cyrillic, hence the transliteration of "yo" to "e"). IUPAC adopted the name on 30 May 2012. The name and symbol had previously been proposed for element 102 (nobelium) but were not accepted by IUPAC at that time.
It is a transactinide in the p-block of the periodic table. It is in period 7 and is the heaviest known member of the carbon group. Initial chemical studies in 2007–2008 indicated that flerovium was unexpectedly volatile for a group 14 element. More recent results show that flerovium's reaction with gold is similar to that of copernicium, showing it is very volatile and may even be gaseous at standard temperature and pressure. Nonetheless it also seems to show some metallic properties, consistent with it being the heavier homologue of lead.
Very little is known about flerovium, as it can only be produced one atom at a time, either through direct synthesis or through radioactive decay of even heavier elements, and all known isotopes are short-lived. Six isotopes of flerovium are known, ranging in mass number between 284 and 289; the most stable of these, , has a half-life of ~2.1 seconds, but the unconfirmed may have a longer half-life of 19 seconds, which would be one of the longest half-lives of any nuclide in these farthest reaches of the periodic table. Flerovium is predicted to be near the centre of the theorized island of stability, and it is expected that heavier flerovium isotopes, especially the possibly magic , may have even longer half-lives.
Introduction
History
Pre-discovery
In the late 1940s to early 1960s, the early days of making heavier and heavier transuranic elements, it was predicted that since such elements did not occur naturally, they would have shorter and shorter spontaneous fission half-lives, until they stopped existing altogether around element 108 (now called hassium). Initial work in synthesizing the heavier actinides seemed to confirm this. But the nuclear shell model, introduced in 1949 and extensively developed in the late 1960s by William Myers and Władysław Świątecki, stated that protons and neutrons form shells within a nucleus, analogous to electron shells. Noble gases are unreactive due to a full electron shell; similarly, it was theorized that elements with full nuclear shells – those having "magic" numbers of protons or neutrons – would be stabilized against decay. A doubly magic isotope, with magic numbers of both protons and neutrons, would be especially stabilized. Heiner Meldner calculated in 1965 that the next doubly magic isotope after was with 114 protons and 184 neutrons, which would be the centre of an "island of stability". This island of stability, supposedly from copernicium (Z = 112) to oganesson (Z = 118), would come after a long "sea of instability" from mendelevium (Z = 101) to roentgenium (Z = 111), and the flerovium isotopes in it were speculated in 1966 to have half-lives over 108 years. These early predictions fascinated researchers, and led to the first attempt to make flerovium, in 1968 with the reaction . No flerovium atoms were detected; this was thought to be because the compound nucleus only has 174 neutrons instead of the supposed magic 184, and this would have significant impact on the reaction cross section (yield) and half-lives of nuclei produced. It was then 30 more years before flerovium was first made. Later work suggests the islands of stability around hassium and flerovium occur because these nuclei are respectively deformed and oblate, which make them resistant to spontaneous fission, and that the true island of stability for spherical nuclei occurs at around unbibium-306 (122 protons, 184 neutrons).
In the 1970s and 1980s, theoretical studies debated whether element 114 would be a more volatile metal like lead, or an inert gas.
First signs
The first sign of flerovium was found in December 1998 by a team of scientists at Joint Institute for Nuclear Research (JINR), Dubna, Russia, led by Yuri Oganessian, who bombarded a target of plutonium-244 with accelerated nuclei of calcium-48:
+ → * → + 2
This reaction had been tried before, without success; for this 1998 attempt, JINR had upgraded all of its equipment to detect and separate the produced atoms better and bombard the target more intensely. One atom of flerovium, alpha decaying with lifetime 30.4 s, was detected. The decay energy measured was 9.71 MeV, giving an expected half-life of 2–23 s. This observation was assigned to and was published in January 1999. The experiment was later repeated, but an isotope with these decay properties was never observed again, so the exact identity of this activity is unknown. It may have been due to the isomer , but because the presence of a whole series of longer-lived isomers in its decay chain would be rather doubtful, the most likely assignment of this chain is to the 2n channel leading to and electron capture to . This fits well with the systematics and trends of flerovium isotopes, and is consistent with the low beam energy chosen for that experiment, though further confirmation would be desirable via synthesis of in a 248Cm(48Ca,2n) reaction, which would alpha decay to . The RIKEN team reported possible synthesis of isotopes and in 2016 in a 248Cm(48Ca,2n) reaction, but the alpha decay of was missed, alpha decay of to was observed instead of electron capture to , and the assignment to instead of was not certain.
Glenn T. Seaborg, a scientist at Lawrence Berkeley National Laboratory who had been involved in work to make such superheavy elements, had said in December 1997 that "one of his longest-lasting and most cherished dreams was to see one of these magic elements"; he was told of the synthesis of flerovium by his colleague Albert Ghiorso soon after its publication in 1999. Ghiorso later recalled:
Seaborg died two months later, on 25 February 1999.
In March 1999, the same team replaced the target with to make other flerovium isotopes. Two atoms of flerovium were produced as a result, each alpha-decaying with a half-life of 5.5 s. They were assigned as . This activity has not been seen again either, and it is unclear what nucleus was produced. It is possible that it was an isomer 287mFl or from electron capture by 287Fl, leading to 287Nh and 283Rg.
Confirmed discovery
The now-confirmed discovery of flerovium was made in June 1999 when the Dubna team repeated the first reaction from 1998. This time, two atoms of flerovium were produced; they alpha decayed with half-life 2.6 s, different from the 1998 result. This activity was initially assigned to 288Fl in error, due to the confusion regarding the previous observations that were assumed to come from 289Fl. Further work in December 2002 finally allowed a positive reassignment of the June 1999 atoms to 289Fl.
In May 2009, the Joint Working Party (JWP) of IUPAC published a report on the discovery of copernicium in which they acknowledged discovery of the isotope 283Cn. This implied the discovery of flerovium, from the acknowledgement of the data for the synthesis of 287Fl and 291Lv, which decay to 283Cn. The discovery of flerovium-286 and -287 was confirmed in January 2009 at Berkeley. This was followed by confirmation of flerovium-288 and -289 in July 2009 at Gesellschaft für Schwerionenforschung (GSI) in Germany. In 2011, IUPAC evaluated the Dubna team's 1999–2007 experiments. They found the early data inconclusive, but accepted the results of 2004–2007 as flerovium, and the element was officially recognized as having been discovered.
Isotopes
While the method of chemical characterization of a daughter was successful for flerovium and livermorium, and the simpler structure of even–even nuclei made confirmation of oganesson (Z = 118) straightforward, there have been difficulties in establishing the congruence of decay chains from isotopes with odd protons, odd neutrons, or both. To get around this problem with hot fusion, the decay chains from which terminate in spontaneous fission instead of connecting to known nuclei as cold fusion allows, experiments were done in Dubna in 2015 to produce lighter isotopes of flerovium by reaction of 48Ca with 239Pu and 240Pu, particularly 283Fl, 284Fl, and 285Fl; the last had previously been characterized in the 242Pu(48Ca,5n)285Fl reaction at Lawrence Berkeley National Laboratory in 2010. 285Fl was more clearly characterized, while the new isotope 284Fl was found to undergo immediate spontaneous fission, and 283Fl was not observed. This lightest isotope may yet conceivably be produced in the cold fusion reaction 208Pb(76Ge,n)283Fl, which the team at RIKEN in Japan at one point considered investigating: this reaction is expected to have a higher cross-section of 200 fb than the "world record" low of 30 fb for 209Bi(70Zn,n)278Nh, the reaction which RIKEN used for the official discovery of element 113 (nihonium). Alternatively, it might be produced in future as a great-granddaughter of 295120, reachable in the 249Cf(50Ti,4n) reaction. The reaction 239Pu+48Ca has also been suggested as a means to produce 282Fl and 283Fl in the 5n and 4n channels respectively, but so far only the 3n channel leading to 284Fl has been observed.
The Dubna team repeated their investigation of the 240Pu+48Ca reaction in 2017, observing three new consistent decay chains of 285Fl, another decay chain from this nuclide that may pass through some isomeric states in its daughters, a chain that could be assigned to 287Fl (likely from 242Pu impurities in the target), and some spontaneous fissions of which some could be from 284Fl, though other interpretations including side reactions involving evaporation of charged particles are also possible. The alpha decay of 284Fl to spontaneously fissioning 280Cn was finally observed by the Dubna team in 2024.
Naming
Per Mendeleev's nomenclature for unnamed and undiscovered elements, flerovium is sometimes called eka-lead. In 1979, IUPAC published recommendations according to which the element was to be called ununquadium (symbol Uuq), a systematic element name as a placeholder, until the discovery of the element is confirmed and a permanent name is decided on. Most scientists in the field called it "element 114", with the symbol of E114, (114) or 114.
Per IUPAC recommendations, the discoverer(s) of a new element has the right to suggest a name.
After IUPAC recognized the discovery of flerovium and livermorium on 1 June 2011, IUPAC asked the discovery team at JINR to suggest permanent names for the two elements. The Dubna team chose the name flerovium (symbol Fl), after Russia's Flerov Laboratory of Nuclear Reactions (FLNR), named after Soviet physicist Georgy Flyorov (also spelled Flerov); earlier reports claim the element name was directly proposed to honour Flyorov. In accordance with the proposal received from the discoverers, IUPAC officially named flerovium after Flerov Laboratory of Nuclear Reactions, not after Flyorov himself. Flyorov is known for writing to Joseph Stalin in April 1942 and pointing out the silence in scientific journals in the field of nuclear fission in the United States, Great Britain, and Germany. Flyorov deduced that this research must have become classified information in those countries. Flyorov's work and urgings led to the development of the USSR's own atomic bomb project. Flyorov is also known for the discovery of spontaneous fission with Konstantin Petrzhak. The naming ceremony for flerovium and livermorium was held on 24 October 2012 in Moscow.
In a 2015 interview with Oganessian, the host, in preparation to ask a question, said, "You said you had dreamed to name [an element] after your teacher Georgy Flyorov." Without letting the host finish, Oganessian repeatedly said, "I did."
Predicted properties
Very few properties of flerovium or its compounds have been measured; due to its extremely limited and expensive production and the fact that it decays very quickly. A few singular properties have been measured, but for the most part, properties of flerovium remain unknown and only predictions are available.
Nuclear stability and isotopes
The basis of the chemical periodicity in the periodic table is the electron shell closure at each noble gas (atomic numbers 2, 10, 18, 36, 54, 86, and 118): as any further electrons must enter a new shell with higher energy, closed-shell electron configurations are markedly more stable, hence the inertness of noble gases. Protons and neutrons are also known to form closed nuclear shells, so the same happens at nucleon shell closures, which happen at specific nucleon numbers often dubbed "magic numbers". The known magic numbers are 2, 8, 20, 28, 50, and 82 for protons and neutrons; also 126 for neutrons. Nuclei with magic proton and neutron numbers, such as helium-4, oxygen-16, calcium-48, and lead-208, are "doubly magic" and are very stable. This stability is very important for superheavy elements: with no stabilization, half-lives would be expected by exponential extrapolation to be nanoseconds at darmstadtium (element 110), because the ever-increasing electrostatic repulsion between protons overcomes the limited-range strong nuclear force that holds nuclei together. The next closed nucleon shells (magic numbers) are thought to denote the centre of the long-sought island of stability, where half-lives to alpha decay and spontaneous fission lengthen again.
Initially, by analogy with neutron magic number 126, the next proton shell was also expected at element 126, too far beyond the synthesis capabilities of the mid-20th century to get much theoretical attention. In 1966, new values for the potential and spin–orbit interaction in this region of the periodic table contradicted this and predicted that the next proton shell would instead be at element 114, and that nuclei in this region would be relatively stable against spontaneous fission. The expected closed neutron shells in this region were at neutron number 184 or 196, making 298Fl and 310Fl candidates for being doubly magic. 1972 estimates predicted a half-life of around 1 year for 298Fl, which was expected to be near an island of stability centered near 294Ds (with a half-life around 1010 years, comparable to 232Th). After making the first isotopes of elements 112–118 at the turn of the 21st century, it was found that these neutron-deficient isotopes were stabilized against fission. In 2008 it was thus hypothesized that the stabilization against fission of these nuclides was due to their oblate nuclei, and that a region of oblate nuclei was centred on 288Fl. Also, new theoretical models showed that the expected energy gap between the proton orbitals 2f7/2 (filled at element 114) and 2f5/2 (filled at element 120) was smaller than expected, so element 114 no longer appeared to be a stable spherical closed nuclear shell. The next doubly magic nucleus is now expected to be around 306Ubb, but this nuclide's expected short half-life and low production cross section make its synthesis challenging. Still, the island of stability is expected to exist in this region, and nearer its centre (which has not been approached closely enough yet) some nuclides, such as 291Mc and its alpha- and beta-decay daughters, may be found to decay by positron emission or electron capture and thus move into the centre of the island. Due to the expected high fission barriers, any nucleus in this island of stability would decay exclusively by alpha decay and perhaps some electron capture and beta decay, both of which would bring the nuclei closer to the beta-stability line where the island is expected to be. Electron capture is needed to reach the island, which is problematic because it is not certain that electron capture is a major decay mode in this region of the chart of nuclides.
Experiments were done in 2000–2004 at Flerov Laboratory of Nuclear Reactions in Dubna studying the fission properties of the compound nucleus 292Fl by bombarding 244Pu with accelerated 48Ca ions. A compound nucleus is a loose combination of nucleons that have not yet arranged themselves into nuclear shells. It has no internal structure and is held together only by the collision forces between the two nuclei. Results showed how such nuclei fission mainly by expelling doubly magic or nearly doubly magic fragments such as 40Ca, 132Sn, 208Pb, or 209Bi. It was also found that 48Ca and 58Fe projectiles had a similar yield for the fusion-fission pathway, suggesting possible future use of 58Fe projectiles in making superheavy elements. It has also been suggested that a neutron-rich flerovium isotope can be formed by quasifission (partial fusion followed by fission) of a massive nucleus. Recently it has been shown that multi-nucleon transfer reactions in collisions of actinide nuclei (such as uranium and curium) might be used to make neutron-rich superheavy nuclei in the island of stability, though production of neutron-rich nobelium or seaborgium is more likely.
Theoretical estimates of alpha decay half-lives of flerovium isotopes, support the experimental data.
The fission-survived isotope 298Fl, long expected to be doubly magic, is predicted to have alpha decay half-life ~17 days. Making 298Fl directly by a fusion–evaporation pathway is currently impossible: no known combination of target and stable projectile can give 184 neutrons for the compound nucleus, and radioactive projectiles such as 50Ca (half-life 14 s) cannot yet be used in the needed quantity and intensity. One possibility for making the theorized long-lived nuclei of copernicium (291Cn and 293Cn) and flerovium near the middle of the island, is using even heavier targets such as 250Cm, 249Bk, 251Cf, and 254Es, that when fused with 48Ca would yield isotopes such as 291Mc and 291Fl (as decay products of 299Uue, 295Ts, and 295Lv), which may have just enough neutrons to alpha decay to nuclides close enough to the centre of the island to possibly undergo electron capture and move inward to the centre. However, reaction cross sections would be small and little is yet known about the decay properties of superheavies near the beta-stability line. This may be the current best hope to synthesize nuclei in the island of stability, but it is speculative and may or may not work in practice. Another possibility is to use controlled nuclear explosions to get the high neutron flux needed to make macroscopic amounts of such isotopes. This would mimic the r-process where the actinides were first produced in nature and the gap of instability after polonium bypassed, as it would bypass the gaps of instability at 258–260Fm and at mass number 275 (atomic numbers 104 to 108). Some such isotopes (especially 291Cn and 293Cn) may even have been synthesized in nature, but would decay far too quickly (with half-lives of only thousands of years) and be produced in far too small quantities (~10−12 the abundance of lead) to be detectable today outside cosmic rays.
Atomic and physical
Flerovium is in group 14 in the periodic table, below carbon, silicon, germanium, tin, and lead. Every previous group 14 element has 4 electrons in its valence shell, hence valence electron configuration ns2np2. For flerovium, the trend will continue and the valence electron configuration is predicted as 7s27p2; flerovium will be similar to its lighter congeners in many ways. Differences are likely to arise; a large contributor is spin–orbit (SO) interaction—mutual interaction between the electrons' motion and spin. It is especially strong in superheavy elements, because the electrons move faster than in lighter atoms, at speeds comparable to the speed of light. For flerovium, it lowers the 7s and the 7p electron energy levels (stabilizing the corresponding electrons), but two of the 7p electron energy levels are stabilized more than the other four. The stabilization of the 7s electrons is called the inert pair effect, and the effect "tearing" the 7p subshell into the more and less stabilized parts is called subshell splitting. Computational chemists see the split as a change of the second (azimuthal) quantum number from 1 to and for the more stabilized and less stabilized parts of the 7p subshell, respectively. For many theoretical purposes, the valence electron configuration may be represented to reflect the 7p subshell split as 7s7p. These effects cause flerovium's chemistry to be somewhat different from that of its lighter neighbours.
Because the spin–orbit splitting of the 7p subshell is very large in flerovium, and both of flerovium's filled orbitals in the 7th shell are stabilized relativistically; the valence electron configuration of flerovium may be considered to have a completely filled shell. Its first ionization energy of should be the second-highest in group 14. The 6d electron levels are also destabilized, leading to some early speculations that they may be chemically active, though newer work suggests this is unlikely. Because the first ionization energy is higher than in silicon and germanium, though still lower than in carbon, it has been suggested that flerovium could be classed as a metalloid.
Flerovium's closed-shell electron configuration means metallic bonding in metallic flerovium is weaker than in the elements before and after; so flerovium is expected to have a low boiling point, and has recently been suggested to be possibly a gaseous metal, similar to predictions for copernicium, which also has a closed-shell electron configuration. Flerovium's melting and boiling points were predicted in the 1970s to be around 70 and 150 °C, significantly lower than for the lighter group 14 elements (lead has 327 and 1749 °C), and continuing the trend of decreasing boiling points down the group. Earlier studies predicted a boiling point of ~1000 °C or 2840 °C, but this is now considered unlikely because of the expected weak metallic bonding and that group trends would expect flerovium to have low sublimation enthalpy. Preliminary 2021 calculations predicted that flerovium should have melting point −73 °C (lower than mercury at −39 °C and copernicium, predicted 10 ± 11 °C) and boiling point 107 °C, which would make it a liquid metal. Like mercury, radon, and copernicium, but not lead and oganesson (eka-radon), flerovium is calculated to have no electron affinity.
A 2010 study published calculations predicting a hexagonal close-packed crystal structure for flerovium due to spin–orbit coupling effects, and a density of 9.928 g/cm3, though this was noted to be probably slightly too low. Newer calculations published in 2017 expected flerovium to crystallize in face-centred cubic crystal structure like its lighter congener lead, and calculations published in 2022 predicted a density of 11.4 ± 0.3 g/cm3, similar to lead (11.34 g/cm3). These calculations found that the face-centred cubic and hexagonal close-packed structures should have nearly the same energy, a phenomenon reminiscent of the noble gases. These calculations predict that hexagonal close-packed flerovium should be a semiconductor, with a band gap of 0.8 ± 0.3 eV. (Copernicium is also predicted to be a semiconductor.) These calculations predict that the cohesive energy of flerovium should be around −0.5 ± 0.1 eV; this is similar to that predicted for oganesson (−0.45 eV), larger than that predicted for copernicium (−0.38 eV), but smaller than that of mercury (−0.79 eV). The melting point was calculated as 284 ± 50 K (11 ± 50 °C), so that flerovium is probably a liquid at room temperature, although the boiling point was not determined.
The electron of a hydrogen-like flerovium ion (Fl113+; remove all but one electron) is expected to move so fast that its mass is 1.79 times that of a stationary electron, due to relativistic effects. (The figures for hydrogen-like lead and tin are expected to be 1.25 and 1.073 respectively.) Flerovium would form weaker metal–metal bonds than lead and would be adsorbed less on surfaces.
Chemical
Flerovium is the heaviest known member of group 14, below lead, and is projected to be the second member of the 7p series of elements. Nihonium and flerovium are expected to form a very short subperiod corresponding to the filling of the 7p1/2 orbital, coming between the filling of the 6d5/2 and 7p3/2 subshells. Their chemical behaviour is expected to be very distinctive: nihonium's homology to thallium has been called "doubtful" by computational chemists, while flerovium's to lead has been called only "formal".
The first five group 14 members show a +4 oxidation state and the latter members have increasingly prominent +2 chemistry due to onset of the inert pair effect. For tin, the +2 and +4 states are similar in stability, and lead(II) is the most stable of all the chemically well-understood +2 oxidation states in group 14. The 7s orbitals are very highly stabilized in flerovium, so a very large sp3 orbital hybridization is needed to achieve a +4 oxidation state, so flerovium is expected to be even more stable than lead in its strongly predominant +2 oxidation state and its +4 oxidation state should be highly unstable. For example, the dioxide (FlO2) is expected to be highly unstable to decomposition into its constituent elements (and would not be formed by direct reaction of flerovium with oxygen), and flerovane (FlH4), which should have Fl–H bond lengths of 1.787 Å and would be the heaviest homologue of methane (the lighter compounds include silane, germane and stannane), is predicted to be more thermodynamically unstable than plumbane, spontaneously decomposing to flerovium(II) hydride (FlH2) and H2. The tetrafluoride FlF4 would have bonding mostly due to sd hybridizations rather than sp3 hybridizations, and its decomposition to the difluoride and fluorine gas would be exothermic. The other tetrahalides (for example, FlCl4 is destabilized by about 400 kJ/mol) decompose similarly. The corresponding polyfluoride anion should be unstable to hydrolysis in aqueous solution, and flerovium(II) polyhalide anions such as and are predicted to form preferentially in solutions. The sd hybridizations were suggested in early calculations, as flerovium's 7s and 6d electrons share about the same energy, which would allow a volatile hexafluoride to form, but later calculations do not confirm this possibility. In general, spin–orbit contraction of the 7p1/2 orbital should lead to smaller bond lengths and larger bond angles: this has been theoretically confirmed in FlH2. Still, even FlH2 should be relativistically destabilized by 2.6 eV to below Fl+H2; the large spin–orbit effects also break down the usual singlet–triplet divide in the group 14 dihydrides. FlF2 and FlCl2 are predicted to be more stable than FlH2.
Due to relativistic stabilization of flerovium's 7s27p valence electron configuration, the 0 oxidation state should also be more stable for flerovium than for lead, as the 7p1/2 electrons begin to also have a mild inert pair effect: this stabilization of the neutral state may bring about some similarities between the behavior of flerovium and the noble gas radon. Due to flerovium's expected relative inertness, diatomic compounds FlH and FlF should have lower energies of dissociation than the corresponding lead compounds PbH and PbF. Flerovium(IV) should be even more electronegative than lead(IV); lead(IV) has electronegativity 2.33 on the Pauling scale, though the lead(II) value is only 1.87. Flerovium could be a noble metal.
Flerovium(II) should be more stable than lead(II), and halides FlX+, FlX2, , and (X = Cl, Br, I) are expected to form readily. The fluorides would undergo strong hydrolysis in aqueous solution. All flerovium dihalides are expected to be stable; the difluoride being water-soluble. Spin–orbit effects would destabilize the dihydride (FlH2) by almost . In aqueous solution, the oxyanion flerovite () would also form, analogous to plumbite. Flerovium(II) sulfate (FlSO4) and sulfide (FlS) should be very insoluble in water, and flerovium(II) acetate (Fl(C2H3O2)2) and nitrate (Fl(NO3)2) should be quite water-soluble. The standard electrode potential for reduction of Fl2+ ion to metallic flerovium is estimated to be around +0.9 V, confirming the increased stability of flerovium in the neutral state. In general, due to relativistic stabilization of the 7p1/2 spinor, Fl2+ is expected to have properties intermediate between those of Hg2+ or Cd2+ and its lighter congener Pb2+.
Experimental chemistry
Flerovium is currently the last element whose chemistry has been experimentally investigated, though studies so far are not conclusive. Two experiments were done in April–May 2007 in a joint FLNR-PSI collaboration to study copernicium chemistry. The first experiment used the reaction 242Pu(48Ca,3n)287Fl; and the second, 244Pu(48Ca,4n)288Fl: these reactions give short-lived flerovium isotopes whose copernicium daughters would then be studied. Adsorption properties of the resultant atoms on a gold surface were compared to those of radon, as it was then expected that copernicium's full-shell electron configuration would lead to noble-gas like behavior. Noble gases interact with metal surfaces very weakly, which is uncharacteristic of metals.
The first experiment found 3 atoms of 283Cn but seemingly also 1 atom of 287Fl. This was a surprise; transport time for the product atoms is ~2 s, so the flerovium should have decayed to copernicium before adsorption. In the second reaction, 2 atoms of 288Fl and possibly 1 of 289Fl were seen. Two of the three atoms showed adsorption characteristics associated with a volatile, noble-gas-like element, which has been suggested but is not predicted by more recent calculations. These experiments gave independent confirmation for the discovery of copernicium, flerovium, and livermorium via comparison with published decay data. Further experiments in 2008 to confirm this important result detected 1 atom of 289Fl, and supported previous data showing flerovium had a noble-gas-like interaction with gold.
Empirical support for a noble-gas-like flerovium soon weakened. In 2009 and 2010, the FLNR-PSI collaboration synthesized more flerovium to follow up their 2007 and 2008 studies. In particular, the first three flerovium atoms made in the 2010 study suggested again a noble-gas-like character, but the complete set taken together resulted in a more ambiguous interpretation, unusual for a metal in the carbon group but not fully like a noble gas in character. In their paper, the scientists refrained from calling flerovium's chemical properties "close to those of noble gases", as had previously been done in the 2008 study. Flerovium's volatility was again measured through interactions with a gold surface, and provided indications that the volatility of flerovium was comparable to that of mercury, astatine, and the simultaneously investigated copernicium, which had been shown in the study to be a very volatile noble metal, conforming to its being the heaviest known group 12 element. Still, it was pointed out that this volatile behavior was not expected for a usual group 14 metal.
In experiments in 2012 at GSI, flerovium's chemistry was found to be more metallic than noble-gas-like. Jens Volker Kratz and Christoph Düllmann specifically named copernicium and flerovium as being in a new category of "volatile metals"; Kratz even speculated that they might be gases at standard temperature and pressure. These "volatile metals", as a category, were expected to fall between normal metals and noble gases in terms of adsorption properties. Contrary to the 2009 and 2010 results, it was shown in the 2012 experiments that the interactions of flerovium and copernicium respectively with gold were about equal. Further studies showed that flerovium was more reactive than copernicium, in contradiction to previous experiments and predictions.
In a 2014 paper detailing the experimental results of the chemical characterization of flerovium, the GSI group wrote: "[flerovium] is the least reactive element in the group, but still a metal." Nevertheless, in a 2016 conference about chemistry and physics of heavy and superheavy elements, Alexander Yakushev and Robert Eichler, two scientists who had been active at GSI and FLNR in determining flerovium's chemistry, still urged caution based on the inconsistencies of the various experiments previously listed, noting that the question of whether flerovium was a metal or a noble gas was still open with the known evidence: one study suggested a weak noble-gas-like interaction between flerovium and gold, while the other suggested a stronger metallic interaction. The longer-lived isotope has been considered of interest for future radiochemical studies.
Experiments published in 2022 suggest that flerovium is a metal, exhibiting lower reactivity towards gold than mercury, but higher reactivity than radon. The experiments could not identify if the adsorption was due to elemental flerovium (considered more likely), or if it was due to a flerovium compound such as FlO that was more reactive towards gold than elemental flerovium, but both scenarios involve flerovium forming chemical bonds.
| Physical sciences | Group 14 | Chemistry |
77474 | https://en.wikipedia.org/wiki/Nihonium | Nihonium | Nihonium is a synthetic chemical element; it has the symbol Nh and atomic number 113. It is extremely radioactive: its most stable known isotope, nihonium-286, has a half-life of about 10 seconds. In the periodic table, nihonium is a transactinide element in the p-block. It is a member of period 7 and group 13.
Nihonium was first reported to have been created in experiments carried out between 14 July and 10 August 2003, by a Russian–American collaboration at the Joint Institute for Nuclear Research (JINR) in Dubna, Russia, working in collaboration with the Lawrence Livermore National Laboratory in Livermore, California, and on 23 July 2004, by a team of Japanese scientists at Riken in Wakō, Japan. The confirmation of their claims in the ensuing years involved independent teams of scientists working in the United States, Germany, Sweden, and China, as well as the original claimants in Russia and Japan. In 2015, the IUPAC/IUPAP Joint Working Party recognised the element and assigned the priority of the discovery and naming rights for the element to Riken. The Riken team suggested the name nihonium in 2016, which was approved in the same year. The name comes from the common Japanese name for .
Very little is known about nihonium, as it has been made only in very small amounts that decay within seconds. The anomalously long lives of some superheavy nuclides, including some nihonium isotopes, are explained by the island of stability theory. Experiments to date have supported the theory, with the half-lives of the confirmed nihonium isotopes increasing from milliseconds to seconds as neutrons are added and the island is approached. Nihonium has been calculated to have similar properties to its homologues boron, aluminium, gallium, indium, and thallium. All but boron are post-transition metals, and nihonium is expected to be a post-transition metal as well. It should also show several major differences from them; for example, nihonium should be more stable in the +1 oxidation state than the +3 state, like thallium, but in the +1 state nihonium should behave more like silver and astatine than thallium. Preliminary experiments have shown that elemental nihonium is not very volatile, and that it is less reactive than its lighter homologue thallium.
Introduction
History
Early indications
The syntheses of elements 107 to 112 were conducted at the GSI Helmholtz Centre for Heavy Ion Research in Darmstadt, Germany, from 1981 to 1996. These elements were made by cold fusion reactions, in which targets made of lead and bismuth, which are around the stable configuration of 82 protons, are bombarded with heavy ions of period 4 elements. This creates fused nuclei with low excitation energies due to the stability of the targets' nuclei, significantly increasing the yield of superheavy elements. Cold fusion was pioneered by Yuri Oganessian and his team in 1974 at the Joint Institute for Nuclear Research (JINR) in Dubna, Soviet Union. Yields from cold fusion reactions were found to decrease significantly with increasing atomic number; the resulting nuclei were severely neutron-deficient and short-lived. The GSI team attempted to synthesise element 113 via cold fusion in 1998 and 2003, bombarding bismuth-209 with zinc-70; both attempts were unsuccessful.
Faced with this problem, Oganessian and his team at the JINR turned their renewed attention to the older hot fusion technique, in which heavy actinide targets were bombarded with lighter ions. Calcium-48 was suggested as an ideal projectile, because it is very neutron-rich for a light element (combined with the already neutron-rich actinides) and would minimise the neutron deficiencies of the nuclides produced. Being doubly magic, it would confer benefits in stability to the fused nuclei. In collaboration with the team at the Lawrence Livermore National Laboratory (LLNL) in Livermore, California, United States, they made an attempt on element 114 (which was predicted to be a magic number, closing a proton shell, and more stable than element 113).
In 1998, the JINR–LLNL collaboration started their attempt on element 114, bombarding a target of plutonium-244 with ions of calcium-48:
+ → 292114* → 290114 + 2 + e− → 290113 + νe ?
A single atom was observed which was thought to be the isotope 289114: the results were published in January 1999. Despite numerous attempts to repeat this reaction, an isotope with these decay properties has never again been found, and the exact identity of this activity is unknown. A 2016 paper by Sigurd Hofmann et al. considered that the most likely explanation of the 1998 result is that two neutrons were emitted by the produced compound nucleus, leading to 290114 and electron capture to 290113, while more neutrons were emitted in all other produced chains. This would have been the first report of a decay chain from an isotope of element 113, but it was not recognised at the time, and the assignment is still uncertain. A similar long-lived activity observed by the JINR team in March 1999 in the 242Pu + 48Ca reaction may be due to the electron-capture daughter of 287114, 287113; this assignment is also tentative.
JINR–LLNL collaboration
The now-confirmed discovery of element 114 was made in June 1999 when the JINR team repeated the first 244Pu + 48Ca reaction from 1998; following this, the JINR team used the same hot fusion technique to synthesize elements 116 and 118 in 2000 and 2002 respectively via the 248Cm + 48Ca and 249Cf + 48Ca reactions. They then turned their attention to the missing odd-numbered elements, as the odd protons and possibly neutrons would hinder decay by spontaneous fission and result in longer decay chains.
The first report of element 113 was in August 2003, when it was identified as an alpha decay product of element 115. Element 115 had been produced by bombarding a target of americium-243 with calcium-48 projectiles. The JINR–LLNL collaboration published its results in February 2004:
+ → 291115* → 288115 + 3 → 284113 +
+ → 291115* → 287115 + 4 → 283113 +
Four further alpha decays were observed, ending with the spontaneous fission of isotopes of element 105, dubnium.
Riken
While the JINR–LLNL collaboration had been studying fusion reactions with 48Ca, a team of Japanese scientists at the Riken Nishina Center for Accelerator-Based Science in Wakō, Japan, led by Kōsuke Morita had been studying cold fusion reactions. Morita had previously studied the synthesis of superheavy elements at the JINR before starting his own team at Riken. In 2001, his team confirmed the GSI's discoveries of elements 108, 110, 111, and 112. They then made a new attempt on element 113, using the same 209Bi + 70Zn reaction that the GSI had attempted unsuccessfully in 1998. Despite the much lower yield expected than for the JINR's hot fusion technique with calcium-48, the Riken team chose to use cold fusion as the synthesised isotopes would alpha decay to known daughter nuclides and make the discovery much more certain, and would not require the use of radioactive targets. In particular, the isotope 278113 expected to be produced in this reaction would decay to the known 266Bh, which had been synthesised in 2000 by a team at the Lawrence Berkeley National Laboratory (LBNL) in Berkeley.
The bombardment of 209Bi with 70Zn at Riken began in September 2003. The team detected a single atom of 278113 in July 2004 and published their results that September:
+ → 279113* → 278113 +
The Riken team observed four alpha decays from 278113, creating a decay chain passing through 274Rg, 270Mt, and 266Bh before terminating with the spontaneous fission of 262Db. The decay data they observed for the alpha decay of 266Bh matched the 2000 data, lending support for their claim. Spontaneous fission of its daughter 262Db had not been previously known; the American team had observed only alpha decay from this nuclide.
Road to confirmation
When the discovery of a new element is claimed, the Joint Working Party (JWP) of the International Union of Pure and Applied Chemistry (IUPAC) and the International Union of Pure and Applied Physics (IUPAP) assembles to examine the claims according to their criteria for the discovery of a new element, and decides scientific priority and naming rights for the elements. According to the JWP criteria, a discovery must demonstrate that the element has an atomic number different from all previously observed values. It should also preferably be repeated by other laboratories, although this requirement has been waived where the data is of very high quality. Such a demonstration must establish properties, either physical or chemical, of the new element and establish that they are those of a previously unknown element. The main techniques used to demonstrate atomic number are cross-reactions (creating claimed nuclides as parents or daughters of other nuclides produced by a different reaction) and anchoring decay chains to known daughter nuclides. For the JWP, priority in confirmation takes precedence over the date of the original claim. Both teams set out to confirm their results by these methods.
2004–2008
In June 2004 and again in December 2005, the JINR–LLNL collaboration strengthened their claim for the discovery of element 113 by conducting chemical experiments on 268Db, the final decay product of 288115. This was valuable as none of the nuclides in this decay chain were previously known, so that their claim was not supported by any previous experimental data, and chemical experimentation would strengthen the case for their claim, since the chemistry of dubnium is known. 268Db was successfully identified by extracting the final decay products, measuring spontaneous fission (SF) activities and using chemical identification techniques to confirm that they behave like a group 5 element (dubnium is known to be in group 5). Both the half-life and decay mode were confirmed for the proposed 268Db which lends support to the assignment of the parent and daughter nuclei to elements 115 and 113 respectively. Further experiments at the JINR in 2005 confirmed the observed decay data.
In November and December 2004, the Riken team studied the 205Tl + 70Zn reaction, aiming the zinc beam onto a thallium rather than a bismuth target, in an effort to directly produce 274Rg in a cross-bombardment as it is the immediate daughter of 278113. The reaction was unsuccessful, as the thallium target was physically weak compared to the more commonly used lead and bismuth targets, and it deteriorated significantly and became non-uniform in thickness. The reasons for this weakness are unknown, given that thallium has a higher melting point than bismuth. The Riken team then repeated the original 209Bi + 70Zn reaction and produced a second atom of 278113 in April 2005, with a decay chain that again terminated with the spontaneous fission of 262Db. The decay data were slightly different from those of the first chain: this could have been because an alpha particle escaped from the detector without depositing its full energy, or because some of the intermediate decay products were formed in metastable isomeric states.
In 2006, a team at the Heavy Ion Research Facility in Lanzhou, China, investigated the 243Am + 26Mg reaction, producing four atoms of 266Bh. All four chains started with an alpha decay to 262Db; three chains ended there with spontaneous fission, as in the 278113 chains observed at Riken, while the remaining one continued via another alpha decay to 258Lr, as in the 266Bh chains observed at LBNL.
In June 2006, the JINR–LLNL collaboration claimed to have synthesised a new isotope of element 113 directly by bombarding a neptunium-237 target with accelerated calcium-48 nuclei:
+ → 285113* → 282113 + 3
Two atoms of 282113 were detected. The aim of this experiment had been to synthesise the isotopes 281113 and 282113 that would fill in the gap between isotopes produced via hot fusion (283113 and 284113) and cold fusion (278113). After five alpha decays, these nuclides would reach known isotopes of lawrencium, assuming that the decay chains were not terminated prematurely by spontaneous fission. The first decay chain ended in fission after four alpha decays, presumably originating from 266Db or its electron-capture daughter 266Rf. Spontaneous fission was not observed in the second chain even after four alpha decays. A fifth alpha decay in each chain could have been missed, since 266Db can theoretically undergo alpha decay, in which case the first decay chain would have ended at the known 262Lr or 262No and the second might have continued to the known long-lived 258Md, which has a half-life of 51.5 days, longer than the duration of the experiment: this would explain the lack of a spontaneous fission event in this chain. In the absence of direct detection of the long-lived alpha decays, these interpretations remain unconfirmed, and there is still no known link between any superheavy nuclides produced by hot fusion and the well-known main body of the chart of nuclides.
2009–2015
The JWP published its report on elements 113–116 and 118 in 2011. It recognised the JINR–LLNL collaboration as having discovered elements 114 and 116, but did not accept either team's claim to element 113 and did not accept the JINR–LLNL claims to elements 115 and 118. The JINR–LLNL claim to elements 115 and 113 had been founded on chemical identification of their daughter dubnium, but the JWP objected that current theory could not distinguish between superheavy group 4 and group 5 elements by their chemical properties with enough confidence to allow this assignment. The decay properties of all the nuclei in the decay chain of element 115 had not been previously characterised before the JINR experiments, a situation which the JWP generally considers "troublesome, but not necessarily exclusive", and with the small number of atoms produced with neither known daughters nor cross-reactions the JWP considered that their criteria had not been fulfilled. The JWP did not accept the Riken team's claim either due to inconsistencies in the decay data, the small number of atoms of element 113 produced, and the lack of unambiguous anchors to known isotopes.
In early 2009, the Riken team synthesised the decay product 266Bh directly in the 248Cm + 23Na reaction to establish its link with 278113 as a cross-bombardment. They also established the branched decay of 262Db, which sometimes underwent spontaneous fission and sometimes underwent the previously known alpha decay to 258Lr.
In late 2009, the JINR–LLNL collaboration studied the 249Bk + 48Ca reaction in an effort to produce element 117, which would decay to elements 115 and 113 and bolster their claims in a cross-reaction. They were now joined by scientists from Oak Ridge National Laboratory (ORNL) and Vanderbilt University, both in Tennessee, United States, who helped procure the rare and highly radioactive berkelium target necessary to complete the JINR's calcium-48 campaign to synthesise the heaviest elements on the periodic table. Two isotopes of element 117 were synthesised, decaying to element 115 and then element 113:
+ → 297117* → 294117 + 3 → 290115 + α → 286113 + α
+ → 297117* → 293117 + 4 → 289115 + α → 285113 + α
The new isotopes 285113 and 286113 produced did not overlap with the previously claimed 282113, 283113, and 284113, so this reaction could not be used as a cross-bombardment to confirm the 2003 or 2006 claims.
In March 2010, the Riken team again attempted to synthesise 274Rg directly through the 205Tl + 70Zn reaction with upgraded equipment; they failed again and abandoned this cross-bombardment route.
After 450 more days of irradiation of bismuth with zinc projectiles, Riken produced and identified another 278113 atom in August 2012. Although electricity prices had soared since the 2011 Tōhoku earthquake and tsunami, and Riken had ordered the shutdown of the accelerator programs to save money, Morita's team was permitted to continue with one experiment, and they chose their attempt to confirm their synthesis of element 113. In this case, a series of six alpha decays was observed, leading to an isotope of mendelevium:
278113 → + → + → + → + → + → +
This decay chain differed from the previous observations at Riken mainly in the decay mode of 262Db, which was previously observed to undergo spontaneous fission, but in this case instead alpha decayed; the alpha decay of 262Db to 258Lr is well-known. The team calculated the probability of accidental coincidence to be 10−28, or totally negligible. The resulting 254Md atom then underwent electron capture to 254Fm, which underwent the seventh alpha decay in the chain to the long-lived 250Cf, which has a half-life of around thirteen years.
The 249Bk + 48Ca experiment was repeated at the JINR in 2012 and 2013 with consistent results, and again at the GSI in 2014. In August 2013, a team of researchers at Lund University in Lund, Sweden, and at the GSI announced that they had repeated the 2003 243Am + 48Ca experiment, confirming the findings of the JINR–LLNL collaboration. The same year, the 2003 experiment had been repeated at the JINR, now also creating the isotope 289115 that could serve as a cross-bombardment for confirming their discovery of the element 117 isotope 293117, as well as its daughter 285113 as part of its decay chain. Confirmation of 288115 and its daughters was published by the team at the LBNL in August 2015.
Approval of discoveries
In December 2015, the conclusions of a new JWP report were published by IUPAC in a press release, in which element 113 was awarded to Riken; elements 115, 117, and 118 were awarded to the collaborations involving the JINR. A joint 2016 announcement by IUPAC and IUPAP had been scheduled to coincide with the publication of the JWP reports, but IUPAC alone decided on an early release because the news of Riken being awarded credit for element 113 had been leaked to Japanese newspapers. For the first time in history, a team of Asian physicists would name a new element. The JINR considered the awarding of element 113 to Riken unexpected, citing their own 2003 production of elements 115 and 113, and pointing to the precedents of elements 103, 104, and 105 where IUPAC had awarded joint credit to the JINR and LBNL. They stated that they respected IUPAC's decision, but reserved determination of their position for the official publication of the JWP reports.
The full JWP reports were published on 21 January 2016. The JWP recognised the discovery of element 113, assigning priority to Riken. They noted that while the individual decay energies of each nuclide in the decay chain of 278113 were inconsistent, their sum was now confirmed to be consistent, strongly suggesting that the initial and final states in 278113 and its daughter 262Db were the same for all three events. The decay of 262Db to 258Lr and 254Md was previously known, firmly anchoring the decay chain of 278113 to known regions of the chart of nuclides. The JWP considered that the JINR–LLNL collaborations of 2004 and 2007, producing element 113 as the daughter of element 115, did not meet the discovery criteria as they had not convincingly determined the atomic numbers of their nuclides through cross-bombardments, which were considered necessary since their decay chains were not anchored to previously known nuclides. They also considered that the previous JWP's concerns over their chemical identification of the dubnium daughter had not been adequately addressed. The JWP recognised the JINR–LLNL–ORNL–Vanderbilt collaboration of 2010 as having discovered elements 117 and 115, and accepted that element 113 had been produced as their daughter, but did not give this work shared credit.
After the publication of the JWP reports, Sergey Dimitriev, the lab director of the Flerov lab at the JINR where the discoveries were made, remarked that he was happy with IUPAC's decision, mentioning the time Riken spent on their experiment and their good relations with Morita, who had learnt the basics of synthesising superheavy elements at the JINR.
The sum argument advanced by the JWP in the approval of the discovery of element 113 was later criticised in a May 2016 study from Lund University and the GSI, as it is only valid if no gamma decay or internal conversion takes place along the decay chain, which is not likely for odd nuclei, and the uncertainty of the alpha decay energies measured in the 278113 decay chain was not small enough to rule out this possibility. If this is the case, similarity in lifetimes of intermediate daughters becomes a meaningless argument, as different isomers of the same nuclide can have different half-lives: for example, the ground state of 180Ta has a half-life of hours, but an excited state 180mTa has never been observed to decay. This study found reason to doubt and criticise the IUPAC approval of the discoveries of elements 115 and 117, but the data from Riken for element 113 was found to be congruent, and the data from the JINR team for elements 115 and 113 to probably be so, thus endorsing the IUPAC approval of the discovery of element 113. Two members of the JINR team published a journal article rebutting these criticisms against the congruence of their data on elements 113, 115, and 117 in June 2017.
Naming
Using Mendeleev's nomenclature for unnamed and undiscovered elements, nihonium would be known as eka-thallium. In 1979, IUPAC published recommendations according to which the element was to be called ununtrium (with the corresponding symbol of Uut), a systematic element name as a placeholder, until the discovery of the element is confirmed and a name is decided on. The recommendations were widely used in the chemical community on all levels, from chemistry classrooms to advanced textbooks, but were mostly ignored among scientists in the field, who called it "element 113", with the symbol of E113, (113), or even simply 113.
Before the JWP recognition of their priority, the Japanese team had unofficially suggested various names: japonium, after their home country; nishinanium, after Japanese physicist Yoshio Nishina, the "founding father of modern physics research in Japan"; and rikenium, after the institute. After the recognition, the Riken team gathered in February 2016 to decide on a name. Morita expressed his desire for the name to honour the fact that element 113 had been discovered in Japan. Japonium was considered, making the connection to Japan easy to identify for non-Japanese, but it was rejected as Jap is considered an ethnic slur. The name nihonium was chosen after an hour of deliberation: it comes from , one of the two Japanese pronunciations for the name of Japan. The discoverers also intended to reference the support of their research by the Japanese people (Riken being almost entirely government-funded), recover lost pride and trust in science among those who were affected by the Fukushima Daiichi nuclear disaster, and honour Japanese chemist Masataka Ogawa's 1908 discovery of rhenium, which he named "nipponium" with symbol Np after the other Japanese pronunciation of Japan's name. As Ogawa's claim had not been accepted, the name "nipponium" could not be reused for a new element, and its symbol Np had since been used for neptunium. In March 2016, Morita proposed the name "nihonium" to IUPAC, with the symbol Nh. The naming realised what had been a national dream in Japanese science ever since Ogawa's claim.
The former president of IUPAP, Cecilia Jarlskog, complained at the Nobel Symposium on Superheavy Elements in Bäckaskog Castle, Sweden, in June 2016 about the lack of openness involved in the process of approving new elements, and stated that she believed that the JWP's work was flawed and should be redone by a new JWP. A survey of physicists determined that many felt that the Lund–GSI 2016 criticisms of the JWP report were well-founded, but it was also generally thought that the conclusions would hold up if the work was redone. Thus the new president, Bruce McKellar, ruled that the proposed names should be released in a joint IUPAP–IUPAC press release. IUPAC and IUPAP publicised the proposal of nihonium that June, and set a five-month term to collect comments, after which the name would be formally established at a conference. The name was officially approved on 28 November 2016. The naming ceremony for the new element was held in Tokyo, Japan, on 14 March 2017, with Naruhito, then the Crown Prince of Japan, in attendance.
Isotopes
Nihonium has no stable or naturally occurring isotopes. Several radioactive isotopes have been synthesised in the laboratory, either by fusing two atoms or by observing the decay of heavier elements. Eight different isotopes of nihonium have been reported with atomic masses 278, 282–287, and 290 (287Nh and 290Nh are unconfirmed); they all decay through alpha decay to isotopes of roentgenium. There have been indications that nihonium-284 can also decay by electron capture to copernicium-284, though estimates of the partial half-life for this branch vary strongly by model. A spontaneous fission branch of nihonium-285 has also been reported.
Stability and half-lives
The stability of nuclei quickly decreases with the increase in atomic number after curium, element 96, whose half-life is over ten thousand times longer than that of any subsequent element. All isotopes with an atomic number above 101 undergo radioactive decay with half-lives of less than 30 hours: this is because of the ever-increasing Coulomb repulsion of protons, so that the strong nuclear force cannot hold the nucleus together against spontaneous fission for long. Calculations suggest that in the absence of other stabilising factors, elements with more than 103 protons should not exist. Researchers in the 1960s suggested that the closed nuclear shells around 114 protons and 184 neutrons should counteract this instability, and create an "island of stability" containing nuclides with half-lives reaching thousands or millions of years. The existence of the island is still unproven, but the existence of the superheavy elements (including nihonium) confirms that the stabilising effect is real, and in general the known superheavy nuclides become longer-lived as they approach the predicted location of the island.
All nihonium isotopes are unstable and radioactive; the heavier nihonium isotopes are more stable than the lighter ones, as they are closer to the centre of the island. The most stable known nihonium isotope, 286Nh, is also the heaviest; it has a half-life of 8 seconds. The isotope 285Nh, as well as the unconfirmed 287Nh and 290Nh, have also been reported to have half-lives of over a second. The isotopes 284Nh and 283Nh have half-lives of 0.90 and 0.12 seconds respectively. The remaining two isotopes have half-lives between 0.1 and 100 milliseconds: 282Nh has a half-life of 61 milliseconds, and 278Nh, the lightest known nihonium isotope, is also the shortest-lived, with a half-life of 2.0 milliseconds. This rapid increase in the half-lives near the closed neutron shell at N = 184 is seen in roentgenium, copernicium, and nihonium (elements 111 through 113), where each extra neutron so far multiplies the half-life by a factor of 5 to 20.
The unknown isotopes in the gap between 278Nh and 282Nh are too heavy to be produced by cold fusion and too light to be produced by hot fusion. The missing 280Nh and 281Nh may be populated as daughters of 284Mc and 285Mc, producible in the 241Am+48Ca reaction, but this has not yet been attempted. Of particular interest is 281Nh, as it is the expected great-granddaughter of 293119, a possible product of the 243Am+54Cr reaction. Production of 282Mc and 283Mc is possible in the 243Am+44Ca reaction (though it has a lower cross-section), and their daughters would be 278Nh (known) and 279Nh. The heavier isotopes 287Nh through 290Nh might be synthesised using charged-particle evaporation, using the 242Pu+48Ca and 244Pu+48Ca reactions where one proton and some neutrons are evaporated.
Predicted properties
Very few properties of nihonium or its compounds have been measured; this is due to its extremely limited and expensive production and the fact it decays very quickly. Properties of nihonium mostly remain unknown and only predictions are available.
Physical and atomic
Nihonium is the first member of the 7p series of elements and the heaviest group 13 element on the periodic table, below boron, aluminium, gallium, indium, and thallium. All the group 13 elements except boron are metals, and nihonium is expected to follow suit. Nihonium is predicted to show many differences from its lighter homologues. The major reason for this is the spin–orbit (SO) interaction, which is especially strong for the superheavy elements, because their electrons move much faster than in lighter atoms, at velocities close to the speed of light. In relation to nihonium atoms, it lowers the 7s and the 7p electron energy levels (stabilising those electrons), but two of the 7p electron energy levels are stabilised more than the other four. The stabilisation of the 7s electrons is called the inert pair effect, and the separation of the 7p subshell into the more and less stabilised parts is called subshell splitting. Computational chemists see the split as a change of the second, azimuthal quantum number l, from 1 to 1/2 and 3/2 for the more and less stabilised parts of the 7p subshell, respectively. The quantum number corresponds to the letter in the electron orbital name: 0 to s, 1 to p, 2 to d, etc. For theoretical purposes, the valence electron configuration may be represented to reflect the 7p subshell split as 7s2 7p1/21. The first ionisation energy of nihonium is expected to be 7.306 eV, the highest among the metals of group 13. Similar subshell splitting should exist for the 6d electron levels, with four being 6d3/2 and six being 6d5/2. Both these levels are raised to be close in energy to the 7s ones, high enough to possibly be chemically active. This would allow for the possibility of exotic nihonium compounds without lighter group 13 analogues.
Periodic trends would predict nihonium to have an atomic radius larger than that of thallium due to it being one period further down the periodic table, but calculations suggest nihonium has an atomic radius of about 170 pm, the same as that of thallium, due to the relativistic stabilisation and contraction of its 7s and 7p1/2 orbitals. Thus, nihonium is expected to be much denser than thallium, with a predicted density of about 16 to 18 g/cm3 compared to thallium's 11.85 g/cm3, since nihonium atoms are heavier than thallium atoms but have the same volume. Bulk nihonium is expected to have a hexagonal close-packed crystal structure, like thallium. The melting and boiling points of nihonium have been predicted to be 430 °C and 1100 °C respectively, exceeding the values for indium and thallium, following periodic trends. Nihonium should have a bulk modulus of 20.8 GPa, about half that of thallium (43 GPa).
Chemical
The chemistry of nihonium is expected to be very different from that of thallium. This difference stems from the spin–orbit splitting of the 7p shell, which results in nihonium being between two relatively inert closed-shell elements (copernicium and flerovium). Nihonium is expected to be less reactive than thallium, because of the greater stabilisation and resultant chemical inactivity of the 7s subshell in nihonium compared to the 6s subshell in thallium. The standard electrode potential for the Nh+/Nh couple is predicted to be 0.6 V. Nihonium should be a rather noble metal.
The metallic group 13 elements are typically found in two oxidation states: +1 and +3. The former results from the involvement of only the single p electron in bonding, and the latter results in the involvement of all three valence electrons, two in the s-subshell and one in the p-subshell. Going down the group, bond energies decrease and the +3 state becomes less stable, as the energy released in forming two additional bonds and attaining the +3 state is not always enough to outweigh the energy needed to involve the s-electrons. Hence, for aluminium and gallium +3 is the most stable state, but +1 gains importance for indium and by thallium it becomes more stable than the +3 state. Nihonium is expected to continue this trend and have +1 as its most stable oxidation state.
The simplest possible nihonium compound is the monohydride, NhH. The bonding is provided by the 7p1/2 electron of nihonium and the 1s electron of hydrogen. The SO interaction causes the binding energy of nihonium monohydride to be reduced by about 1 eV and the nihonium–hydrogen bond length to decrease as the bonding 7p1/2 orbital is relativistically contracted. This is unique among the 7p element monohydrides; all the others have relativistic expansion of the bond length instead of contraction. Another effect of the SO interaction is that the Nh–H bond is expected to have significant pi bonding character (side-on orbital overlap), unlike the almost pure sigma bonding (head-on orbital overlap) in thallium monohydride (TlH). The analogous monofluoride (NhF) should also exist. Nihonium(I) is predicted to be more similar to silver(I) than thallium(I): the Nh+ ion is expected to more willingly bind anions, so that NhCl should be quite soluble in excess hydrochloric acid or ammonia; TlCl is not. In contrast to Tl+, which forms the strongly basic hydroxide (TlOH) in solution, the Nh+ cation should instead hydrolyse all the way to the amphoteric oxide Nh2O, which would be soluble in aqueous ammonia and weakly soluble in water.
The adsorption behaviour of nihonium on gold surfaces in thermochromatographical experiments is expected to be closer to that of astatine than that of thallium. The destabilisation of the 7p3/2 subshell effectively leads to a valence shell closing at the 7s2 7p2 configuration rather than the expected 7s2 7p6 configuration with its stable octet. As such, nihonium, like astatine, can be considered to be one p-electron short of a closed valence shell. Hence, even though nihonium is in group 13, it has several properties similar to the group 17 elements. (Tennessine in group 17 has some group-13-like properties, as it has three valence electrons outside the 7s2 7p2 closed shell.) Nihonium is expected to be able to gain an electron to attain this closed-shell configuration, forming the −1 oxidation state like the halogens (fluorine, chlorine, bromine, iodine, and astatine). This state should be more stable than it is for thallium as the SO splitting of the 7p subshell is greater than that for the 6p subshell. Nihonium should be the most electronegative of the metallic group 13 elements, even more electronegative than tennessine, the period 7 congener of the halogens: in the compound NhTs, the negative charge is expected to be on the nihonium atom rather than the tennessine atom. The −1 oxidation should be more stable for nihonium than for tennessine. The electron affinity of nihonium is calculated to be around 0.68 eV, higher than thallium's at 0.4 eV; tennessine's is expected to be 1.8 eV, the lowest in its group. It is theoretically predicted that nihonium should have an enthalpy of sublimation around 150 kJ/mol and an enthalpy of adsorption on a gold surface around −159 kJ/mol.
Significant 6d involvement is expected in the Nh–Au bond, although it is expected to be more unstable than the Tl–Au bond and entirely due to magnetic interactions. This raises the possibility of some transition metal character for nihonium. On the basis of the small energy gap between the 6d and 7s electrons, the higher oxidation states +3 and +5 have been suggested for nihonium. Some simple compounds with nihonium in the +3 oxidation state would be the trihydride (NhH3), trifluoride (NhF3), and trichloride (NhCl3). These molecules are predicted to be T-shaped and not trigonal planar as their boron analogues are: this is due to the influence of the 6d5/2 electrons on the bonding. The heavier nihonium tribromide (NhBr3) and triiodide (NhI3) are trigonal planar due to the increased steric repulsion between the peripheral atoms; accordingly, they do not show significant 6d involvement in their bonding, though the large 7s–7p energy gap means that they show reduced sp2 hybridisation compared to their boron analogues.
The bonding in the lighter NhX3 molecules can be considered as that of a linear species (similar to HgF2 or ) with an additional Nh–X bond involving the 7p orbital of nihonium perpendicular to the other two ligands. These compounds are all expected to be highly unstable towards the loss of an X2 molecule and reduction to nihonium(I):
NhX3 → NhX + X2
Nihonium thus continues the trend down group 13 of reduced stability of the +3 oxidation state, as all five of these compounds have lower reaction energies than the unknown thallium(III) iodide. The +3 state is stabilised for thallium in anionic complexes such as , and the presence of a possible vacant coordination site on the lighter T-shaped nihonium trihalides is expected to allow a similar stabilisation of and perhaps .
The +5 oxidation state is unknown for all lighter group 13 elements: calculations predict that nihonium pentahydride (NhH5) and pentafluoride (NhF5) should have a square pyramidal molecular geometry, but also that both would be highly thermodynamically unstable to loss of an X2 molecule and reduction to nihonium(III). Again, some stabilisation is expected for anionic complexes, such as . The structures of the nihonium trifluoride and pentafluoride molecules are the same as those for chlorine trifluoride and pentafluoride.
Experimental chemistry
The isotopes 284Nh, 285Nh, and 286Nh have half-lives long enough for chemical investigation. From 2010 to 2012, some preliminary chemical experiments were performed at the JINR to determine the volatility of nihonium. The isotope 284Nh was investigated, made as the daughter of 288Mc produced in the 243Am+48Ca reaction. The nihonium atoms were synthesised in a recoil chamber and then carried along polytetrafluoroethylene (PTFE) capillaries at 70 °C by a carrier gas to the gold-covered detectors. About ten to twenty atoms of 284Nh were produced, but none of these atoms were registered by the detectors, suggesting either that nihonium was similar in volatility to the noble gases (and thus diffused away too quickly to be detected) or, more plausibly, that pure nihonium was not very volatile and thus could not efficiently pass through the PTFE capillaries. Formation of the hydroxide NhOH should ease the transport, as nihonium hydroxide is expected to be more volatile than elemental nihonium, and this reaction could be facilitated by adding more water vapour into the carrier gas. It seems likely that this formation is not kinetically favoured, so the longer-lived isotopes 285Nh and 286Nh were considered more desirable for future experiments.
A 2017 experiment at the JINR, producing 284Nh and 285Nh via the 243Am+48Ca reaction as the daughters of 288Mc and 289Mc, avoided this problem by removing the quartz surface, using only PTFE. No nihonium atoms were observed after chemical separation, implying an unexpectedly large retention of nihonium atoms on PTFE surfaces. This experimental result for the interaction limit of nihonium atoms with a PTFE surface disagrees significantly with previous theory, which expected a lower value of 14.00 kJ/mol. This suggests that the nihonium species involved in the previous experiment was likely not elemental nihonium but rather nihonium hydroxide, and that high-temperature techniques such as vacuum chromatography would be necessary to further probe the behaviour of elemental nihonium. Bromine saturated with boron tribromide has been suggested as a carrier gas for experiments on nihonium chemistry; this oxidises nihonium's lighter congener thallium to thallium(III), providing an avenue to investigate the oxidation states of nihonium, similar to earlier experiments done on the bromides of group 5 elements, including the superheavy dubnium.
A 2024 experiment at the GSI, producing 284Nh via the 243Am+48Ca reaction as daughter of 288Mc, studied the adsorption of nihonium and moscovium on SiO2 and gold surfaces. The adsorption enthalpy of nihonium on SiO2 was determined experimentally as (68% confidence interval). Nihonium was determined to be less reactive with the SiO2 surface than its lighter congener thallium, but more reactive than its closed-shell neighbours copernicium and flerovium. This arises because of the relativistic stabilisation of the 7p1/2 shell.
| Physical sciences | Group 13 | Chemistry |
77510 | https://en.wikipedia.org/wiki/Songbird | Songbird | A songbird is a bird belonging to the suborder Passeri of the perching birds (Passeriformes). Another name that is sometimes seen as the scientific or vernacular name is Oscines, from Latin oscen, "songbird". The Passeriformes contains 5,000 or so species found all over the world, in which the vocal organ typically is developed in such a way as to produce a diverse and elaborate bird song.
Songbirds form one of the two major lineages of extant perching birds (~4,000 species), the other being the Tyranni (~1,000 species), which are most diverse in the Neotropics and absent from many parts of the world. The Tyranni have a simpler syrinx musculature, and while their vocalizations are often just as complex and striking as those of songbirds, they are altogether more mechanical sounding. There is a third perching bird lineage, the Acanthisitti of New Zealand, of which only two species remain alive today. Recent estimates indicate that songbirds originated 50 million years ago. The distribution of their basal lineages suggest that their origin and initial diversification occurred exclusively in the Australian continent and only about 40 million years ago, oscines started to colonize Eurasia, Africa, and eventually the Americas.
Description
The song in this clade is essentially territorial, because it communicates the identity and whereabouts of an individual to other birds, and also signals sexual intentions. Sexual selection among songbirds is highly based on mimetic vocalization. Female preference has shown in some populations to be based on the extent of a male's song repertoire. The larger a male's repertoire, the more females a male individual attracts. It is not to be confused with bird calls that are used for alarms and contact and are especially important in birds that feed or migrate in flocks. While almost all living birds give calls of some sort, well-developed songs are only given by a few lineages outside the songbirds. And still, not all songbirds proffer a call that is distinctly melodious. Songbirds do, however, possess a highly developed vocal organ, the syrinx, that enables their sonorous activity. This organ, also known as a song box, can be found where the windpipe meets diverging bronchial tubes which lead to the lungs. The organ is a solid, bony structure lined with a film of membranes which air passes through as the songbird calls. While the song boxes of songbirds vary in size and intricacy, this does not necessarily determine the songbird's ability to voice their song. Researchers believe this has more to do with the length of the windpipe.
Other birds (especially non-passeriforms) sometimes have songs to attract mates or hold territory, but these are usually simple and repetitive, lacking the variety of many oscine songs. The monotonous repetition of the common cuckoo or little crake can be contrasted with the variety of a nightingale or marsh warbler. However, although many songbirds have songs that are pleasant to the human ear, this is not invariably the case. Many members of the crow family (Corvidae) communicate with croaks or screeches, which sound harsh to humans. Even these, however, have a song of sorts, a softer twitter that is given between courting partners. And even though some parrots (which are not songbirds) can be taught to repeat human speech, vocal mimicry among birds is almost completely restricted to songbirds, some of which (such as the lyrebirds or the aptly named mockingbirds) excel in imitating the sounds of other birds or even environmental noises.
The birds from higher altitudes have evolved thicker downs (also known as jackets) to protect themselves from colder temperatures. Their feathers have outer and inner portions, with the lower down being fluffier and warmer to provide increased warmth.
Song repertoire and courtship
Sexual selection can be broken down into several different studies regarding different aspects of a bird's song. As a result, songs can vary even within a single species. Many believe that song repertoire and cognition have a direct relationship. However, a study published in 2013 has shown that cognitive abilities may not all be directly related to the song repertoire of a songbird. Specifically, spatial learning is said to have an inverse relationship with song repertoire. So for example, this would be an individual who does not migrate as far as others in the species but has a better song repertoire. This suggests an evolutionary trade-off between possible alleles. With natural selection choosing traits best fit for reproductive success, there could be a trade-off in either direction depending on which trait would produce a higher fitness at that time period.
Song repertoire can be attributed to male songbirds as it is one of the main mechanisms of courtship. Song repertoires differ from male individual to male individual and species to species. Some species may typically have large repertoires while others may have significantly smaller ones. Mate choice in female songbirds is a significant realm of study as song abilities are continuously evolving. Males often sing to assert their dominance over other males in competition for a female, sometimes in lieu of a combative episode, and to arouse the female by announcing a readiness to mate. Though less frequent, females have also been known to sing occasionally a duet with a mate as an affirmation of their partnership. While some will sing their song from a familiar perch, other species common to grasslands will sing a familiar song each time they fly. Currently, there have been numerous studies involving songbird repertoires, unfortunately, there has not yet been a concrete evidence to confirm that every songbird species prefers larger repertoires. A conclusion can be made that it can vary between species on whether a larger repertoire is connected to better fitness. With this conclusion, it can be inferred that evolution via natural selection, or sexual selection, favors the ability to retain larger repertoires for these certain species as it leads to higher reproductive success. During times of courtship, it is said that male songbirds increase their repertoire by mimicking other species songs. The better the mimicking ability, retaining ability, and the quantity of other species mimicked has been proven to have a positive relationship with mating success. Female preferences cause the constant improvement of accuracy and presentation of the copied songs. Another theory known as the "song-sharing hypothesis" suggests that females prefer simpler, more homogenous songs that signal a male of familiar territory. As birdsong can be broken into regional dialects through this process of mimicry, the foreign song of a newcomer suggests the lack of territorial possession. This can be costly in the wake of territorial conflicts between disparate songbird populations and may compel a female to prefer a male spouting a familiar song of the area.
Taxonomy and systematics
Sibley and Alquist divided songbirds into two "parvorders", Corvida and Passerida (standard taxonomic practice would rank these as infraorders), distributed in Australo-Papua and Eurasia respectively. Subsequent molecular studies, however, show this treatment to be somewhat erroneous. Passerida is a highly diverse lineage, uniting over one-third of all bird species to include (in 2015) 3,885 species). These are divided into three major superfamilies (though not exactly corresponding to the Sibley-Ahlquist arrangement), in addition to some minor lineages.
In contrast, Sibley & Alquist's "Corvida" is a phylogenetic grade and an artefact of the phenetic methodology. The bulk of the "Corvida" make up the large clade Corvides (812 species as of 2015), which is a sister group to the Passerida. The remaining 15 oscine families (343 species in 2015) form a series of basally branching sister groups to the Corvoid - Passerid clade. All of these groups, which form at least six successively branching basal clades, are found exclusively or predominantly in Australasia. Australian endemics are also prominent among basal lineages in both Corvoids and Passerids, suggesting that songbirds originated and diverged in Australia.
Scrubbirds and lyrebirds, of which there are just two species of each, represent the oldest lineage of songbirds on Earth. The rufous scrubbird, Atrichornis rufescens, is essentially confined to the Gondwana Rainforests of Australia World Heritage Area, occurring in both Queensland and New South Wales sections. It is now only found at elevations above .
One of the earliest known fossil songbirds is Resoviaornis from the Early Oligocene of Poland.
Families
Menuroidea
Menuridae: lyrebirds
Atrichornithidae: scrub birds
Bowerbirds and Australian treecreepers
Climacteridae: Australian treecreepers
Ptilonorhynchidae: bowerbirds
Meliphagoidea: honeyeaters and allies
Maluridae: fairy-wrens, emu-wrens and grasswrens
Meliphagidae: true honeyeaters and chats
Dasyornithidae: bristlebirds
Pardalotidae: pardalotes
Acanthizidae: scrubwrens, thornbills, and gerygones
Australopapuan babblers
Pomatostomidae: Australasian babblers
Logrunners
Orthonychidae: logrunners
Other basal lineages
Cnemophilidae: satinbirds Cnemophilus and Loboparadisea
Melanocharitidae: berrypeckers and longbills
Callaeidae: New Zealand wattlebirds kokako, saddleback and †huia
Notiomystidae: stitchbird
Corvides
Paramythiidae: tit berrypecker and crested berrypeckers
Psophodidae: whipbirds, jewel-babblers and quail-thrushes
Platysteiridae: wattle-eyes and batis
Malaconotidae: bush-shrikes
Machaerirynchidae: boatbills
Vangidae: vangas, woodshrikes, and helmetshrikes
Pityriasidae: Bornean bristlehead
Artamidae: butcherbirds, currawongs and Australian magpie (formerly in Cracticidae)
Rhagologidae: mottled whistler
Aegithinidae: ioras
Campephagidae: cuckooshrikes and trillers
Mohouidae: whiteheads
Neosittidae: sittellas
Eulacestomidae: ploughbill
Oreoicidae: Australo-Papuan bellbirds
Pachycephalidae: whistlers, shrike-thrushes, pitohuis and allies
Laniidae: shrikes
Vireonidae: vireos
Oriolidae: orioles, figbirds and †piopio (formerly Turnagridae)
Dicruridae: drongos
Rhipiduridae: fantails
Monarchidae: monarchs and allies
Platylophidae: jayshrike
Corvidae: crows, magpies, and jays
Corcoracidae: white-winged chough and apostlebird
Melampittidae: melampittas
Ifritidae: ifritabirds
Paradisaeidae: birds of paradise
Passerida
Petroicidae: Australasian robins
Picathartidae: rockfowl
Chaetopidae: rockjumpers
Eupetidae: rail-babbler
Bombycillidae: waxwings and allies
Ptiliogonatidae: silky-flycatchers
Hypocoliidae: hypocolius
Dulidae: palmchat
† Mohoidae: some Hawaiian honeyeaters, Moho and Chaetoptila genera, not closely related to Meliphagidae
Hylocitreidae: hylocitrea
Stenostiridae: fairy-flycatcher and allies
Paridae: tits, chickadees, and titmouse
Remizidae: penduline-tits
Nicatoridae: nicators
Panuridae: bearded reedling
Alaudidae: larks
Pycnonotidae: bulbuls
Hirundinidae: swallows and martins
Pnoepygidae: wren-babblers
Macrosphenidae: crombecs and African warblers
Cettiidae: bush-warblers and allies
Scotocercidae: streaked scrub-warbler
Erythrocercidae: yellow flycatchers
Aegithalidae: long-tailed tits
Phylloscopidae: leaf-warblers and allies. Recently split from Sylviidae.
Acrocephalidae: reed warblers and allies
Locustellidae: grassbirds and allies
Donacobiidae: black-capped donacobius
Bernieridae: Malagasy warblers
Cisticolidae: cisticolas and allies
Timaliidae: babblers
Pellorneidae: ground babblers
Leiothrichidae: laughingthrushes and allies
Sylviidae: Old World warblers
Zosteropidae: white-eyes
Arcanatoridae: dapple-throat and allies
Promeropidae: sugarbirds
Irenidae: fairy-bluebirds
Regulidae: kinglets
Elachuridae: elachuras
Hyliotidae: hyliotas
Troglodytidae: wrens
Polioptilidae: gnatcatchers
Sittidae: nuthatches
Tichodromidae: wallcreeper
Certhiidae: treecreepers
Mimidae: mockingbirds and thrashers
Sturnidae: starlings
Buphagidae: oxpeckers
Turdidae: thrushes and allies
Muscicapidae: Old World flycatchers and chats
Cinclidae: dippers
Chloropseidae: leafbirds
Dicaeidae: flowerpeckers
Nectariniidae: sunbirds
Passeridae: true sparrows
Ploceidae: weavers and widowbirds
Estrildidae: estrildid finches (waxbills, munias, etc.)
Viduidae: indigo birds and whydahs
Peucedramidae: olive warbler
Prunellidae: accentor
Motacillidae: wagtails and pipits
Urocynchramidae: Przevalski's finch
Fringillidae: true finches and Hawaiian honeycreepers (formerly Drepanididae)
Parulidae: New World warblers, for example the black-throated blue warblers and allies
Icteridae: American blackbirds, New World orioles, grackles and cowbirds.
Coerebidae: bananaquit
Emberizidae: buntings
Passerellidae: New World sparrows
Thraupidae: tanagers, true honeycreepers and allies
Calcariidae: snow buntings and longspurs
Cardinalidae: cardinals and allies
| Biology and health sciences | Passerines | null |
77538 | https://en.wikipedia.org/wiki/Perspiration | Perspiration | Perspiration, also known as sweat, is the fluid secreted by sweat glands in the skin of mammals.
Two types of sweat glands can be found in humans: eccrine glands and apocrine glands. The eccrine sweat glands are distributed over much of the body and are responsible for secreting the watery, brackish sweat most often triggered by excessive body temperature. Apocrine sweat glands are restricted to the armpits and a few other areas of the body and produce an odorless, oily, opaque secretion which then gains its characteristic odor from bacterial decomposition.
In humans, sweating is primarily a means of thermoregulation, which is achieved by the water-rich secretion of the eccrine glands. Maximum sweat rates of an adult can be up to per hour or per day, but is less in children prior to puberty. Evaporation of sweat from the skin surface has a cooling effect due to evaporative cooling. Hence, in hot weather, or when the individual's muscles heat up due to exertion, more sweat is produced. Animals with few sweat glands, such as dogs, accomplish similar temperature regulation results by panting, which evaporates water from the moist lining of the oral cavity and pharynx.
Although sweating is found in a wide variety of mammals, relatively few (apart from humans, horses, some primates and some bovidae) produce sweat in order to cool down. In horses, such cooling sweat is created by apocrine glands and contains a wetting agent, the protein latherin which transfers from the skin to the surface of their coats.
Definitions
The words diaphoresis and hidrosis can both mean either perspiration (in which sense they are synonymous with sweating) or excessive perspiration (in which sense they can be either synonymous with hyperhidrosis or differentiable from it only by clinical criteria involved in narrow specialist senses of the words).
Hypohidrosis is decreased sweating from whatever cause.
Focal hyperhidrosis is increased or excessive sweating in certain regions such as the underarm, palms, soles, face, or groin.
Hyperhidrosis is excessive sweating, usually secondary to an underlying condition (in which case it is called secondary hyperhidrosis) and usually involving the body as a whole (in which case it is called generalized hyperhidrosis).
Hidromeiosis is a reduction in sweating that is due to blockages of sweat glands in humid conditions.
A substance or medicine that causes perspiration is a sudorific or sudatory.
Signs and symptoms
Sweat contributes to body odor when it is metabolized by bacteria on the skin. Medications that are used for other treatments and diet also affect odor. Some medical conditions, such as kidney failure and diabetic ketoacidosis, can also affect sweat odor.
Causes
Diaphoresis is a non-specific symptom or sign, which means that it has many possible causes. Some causes of diaphoresis include physical exertion, menopause, fever, ingestion of toxins or irritants, and high environmental temperature. Strong emotions (anger, fear, anxiety) and recall of past trauma can also trigger sweating. This is sometimes referred to as flop sweat.
The vast majority of sweat glands in the body are innervated by sympathetic cholinergic neurons. Sympathetic postganglionic neurons typically secrete norepinephrine and are named sympathetic adrenergic neurons; however, the sympathetic postganglionic neurons that innervate sweat glands secrete acetylcholine and hence are termed sympathetic cholinergic neurons. Sweat glands, piloerector muscles, and some blood vessels are innervated by sympathetic cholinergic neurons.
Pathological sweating and symptoms
Diaphoresis may be associated with some abnormal conditions, such as hyperthyroidism and shock. If it is accompanied by unexplained weight loss, fever/chills, or by palpitations, shortness of breath, unconsciousness, fatigue, dizziness, muscle pain, nausea, vomiting, diarrhea, and chest discomfort, it suggests serious illness.
Diaphoresis is also seen in an acute myocardial infarction (heart attack), from the increased firing of the sympathetic nervous system, and is frequent in serotonin syndrome, which can result in serious sickness or even death. Diaphoresis can also be caused by many types of infections, often accompanied by high fever and/or chills which can trigger the result of hyperthermia. Most infections can cause some degree of diaphoresis and it is a very common symptom in some serious infections such as malaria and tuberculosis. In addition, pneumothorax can cause diaphoresis with splinting of the chest wall.
Neuroleptic malignant syndrome and other malignant diseases (e.g. leukemias) can also cause diaphoresis.
Diabetics relying on insulin shots or oral medications may have low blood sugar (hypoglycemia), which can also cause diaphoresis.
Drugs (including caffeine, morphine, alcohol, antidepressants and certain antipsychotics) may be causes, as well as withdrawal from alcohol, benzodiazepines, nonbenzodiazepines or narcotic painkiller dependencies. Sympathetic nervous system stimulants such as cocaine and amphetamines have also been associated with diaphoresis. Diaphoresis due to ectopic catecholamine is a classic symptom of a pheochromocytoma, a rare tumor of the adrenal gland. Acetylcholinesterase inhibitors (e.g. some insecticides) also cause contraction of sweat gland smooth muscle leading to diaphoresis.
Mercury is well known for its use as a diaphoretic, and was widely used in the 19th and early 20th century by physicians to "purge" the body of an illness. However, due to the high toxicity of mercury, secondary symptoms would manifest, which were erroneously attributed to the former disease that was being treated with mercurials.
Infantile acrodynia (childhood mercury poisoning) is characterized by excessive perspiration. A clinician should immediately consider acrodynia in an afebrile child who is sweating profusely.
Some people can develop a sweat allergy. The allergy is not due to the sweat itself but instead to an allergy-producing protein secreted by bacteria found on the skin. Tannic-acid has been found to suppress the allergic response along with showering.
Hyperhidrosis
Millions of people are affected by hyperhidrosis, but more than half never receive treatment due to embarrassment, lack of awareness, or lack of concern. While it most commonly affects the armpits, feet, and hands, it is possible for someone to experience this condition over their whole body. The face is another common area for hyperhidrosis to be an issue. Sweating uncontrollably is not always expected and may be embarrassing to people with the condition. It can cause both physiological and emotional problems in patients. It is generally inherited. It is not life-threatening, but it is threatening to a person's quality of life. Treatments for hyperhidrosis include antiperspirants, iontophoresis, and surgical removal of sweat glands. In severe cases, botulinum toxin injections or surgical cutting of nerves that stimulate the excessive sweating (endoscopic thoracic sympathectomy) may be an option.
Night sweats
Night sweats, also known as nocturnal hyperhidrosis, is the occurrence of excessive sweating during sleep. The person may or may not also perspire excessively while awake.
One of the most common causes of night sweats in women over 40 is the hormonal changes related to menopause and perimenopause. This is a very common occurrence during the menopausal transition years.
While night sweats might be relatively harmless, it can also be a sign of a serious underlying disease. It is important to distinguish night sweats due to medical causes from those that occur simply because the sleep environment is too warm, either because the bedroom is unusually hot or because there are too many covers on the bed. Night sweats caused by a medical condition or infection can be described as "severe hot flashes occurring at night that can drench sleepwear and sheets, which are not related to the environment". Some of the underlying medical conditions and infections that cause these severe night sweats can be life-threatening and should promptly be investigated by a medical practitioner.
Mechanism
Sweating allows the body to regulate its temperature. Sweating is controlled from a center in the preoptic and anterior regions of the brain's hypothalamus, where thermosensitive neurons are located. The heat-regulatory function of the hypothalamus is also affected by inputs from temperature receptors in the skin. High skin temperature reduces the hypothalamic set point for sweating and increases the gain of the hypothalamic feedback system in response to variations in core temperature. Overall, however, the sweating response to a rise in hypothalamic ('core') temperature is much larger than the response to the same increase in average skin temperature.
Sweating causes a decrease in core temperature through evaporative cooling at the skin surface. As high energy molecules evaporate from the skin, releasing energy absorbed from the body, the skin and superficial vessels decrease in temperature. Cooled venous blood then returns to the body's core and counteracts rising core temperatures.
There are two situations in which the nerves will stimulate the sweat glands, causing perspiration: during physical heat and during emotional stress. In general, emotionally induced sweating is restricted to palms, soles, armpits, and sometimes the forehead, while physical heat-induced sweating occurs throughout the body.
People have an average of two to four million sweat glands, but how much sweat is released by each gland is determined by many factors, including sex, genetics, environmental conditions, age and fitness level. Two of the major contributors to sweat rate are an individual's fitness level and weight. If an individual weighs more, sweat rate is likely to increase because the body must exert more energy to function and there is more body mass to cool down. On the other hand, a fit person will start sweating earlier and more readily. As someone becomes fit, the body becomes more efficient at regulating the body's temperature and sweat glands adapt along with the body's other systems.
Human sweat is not pure water; though it contains no protein, it always contains a small amount (0.2–1%) of solute. When a person moves from a cold climate to a hot climate, adaptive changes occur in the sweating mechanisms of the person. This process is referred to as acclimatization: the maximum rate of sweating increases and its solute composition decreases. The volume of water lost in sweat daily is highly variable, ranging from . The solute loss can be as much as 350mmol/d (or 90mmol/d acclimatised) of sodium under the most extreme conditions. During average intensity exercise, sweat losses can average up to of water/hour. In a cool climate and in the absence of exercise, sodium loss can be very low (less than 5 mmol/d). Sodium concentration in sweat is 30–65 mmol/L, depending on the degree of acclimatisation.
Horses have a thick, waterproofed, hairy coat that would normally block the rapid translocation of sweat water from the skin to the surface of the hair required for evaporative cooling. To solve this, horses have evolved a detergent-like protein, latherin, that they release at high concentrations in their sweat. Their perspiration unlike humans is created by apocrine glands. This protein, by wetting the horses' coat hairs facilitate water flow for cooling evaporation. The presence of this protein can be seen in the lathering that often occurs on the coats of sweating horses, especially when rubbed. In hot conditions, horses during three hours of moderate-intensity exercise can lose of water and of sodium, of chloride and of potassium.
Composition
Sweat is mostly water. A microfluidic model of the eccrine sweat gland provides details on what solutes partition into sweat, their mechanisms of partitioning, and their fluidic transport to the skin surface.
Dissolved in the water are trace amounts of minerals, lactic acid, and urea. Although the mineral content varies, some measured concentrations are: sodium (), potassium (), calcium (), and magnesium ().
Relative to the plasma and extracellular fluid, the concentration of Na+ ions is much lower in sweat (≈40 mM in sweat versus ≈150 mM in plasma and extracellular fluid). Initially, within eccrine glands sweat has a high concentration of Na+ ions. In the sweat ducts, the Na+ ions are re-absorbed into tissue by epithelial sodium channels (ENaC) that are located on the apical membrane of epithelial cells that form the duct (see Fig. 9 of the reference).
Many other trace elements are also excreted in sweat, again an indication of their concentration is (although measurements can vary fifteenfold) zinc (), copper (), iron (), chromium (), nickel (), and lead (). Probably many other less-abundant trace minerals leave the body through sweating with correspondingly lower concentrations. Some exogenous organic compounds make their way into sweat as exemplified by an unidentified odiferous "maple syrup" scented compound in several of the species in the mushroom genus Lactarius. In humans, sweat is hypoosmotic relative to plasma (i.e. less concentrated). Sweat is found at moderately acidic to neutral pH levels, typically between 4.5 and 7.0.
Sweat contains many glycoproteins.
Other functions
Antimicrobial
Sweat may serve an antimicrobial function, like that of earwax or other secretory fluids (e.g., tears, saliva, and milk). It does this through a combination of glycoproteins that either bind directly to, or prevent the binding of microbes to, the skin and seem to form part of the innate immune system.
In 2001, researchers at Eberhard-Karls University in Tübingen, Germany, isolated a large protein called dermcidin from skin. This protein, which could be cleaved into other antimicrobial peptides, was shown to be effective at killing some species of bacteria and fungi that affect humans, including Escherichia coli, Enterococcus faecalis, Staphylococcus aureus, and Candida albicans. It was active at high salt concentrations and in the acidity range of human sweat, where it was present at concentrations of 1–10 mg/ml.
Society and culture
Artificial perspiration
Artificial skin capable of sweating similar to natural sweat rates and with the surface texture and wetting properties of regular skin has been developed for research purposes. Artificial perspiration is also available for in-vitro testing, and contains 19 amino acids and the most abundant minerals and metabolites in sweat.
Diagnostics
There is interest in its use in wearable technology. Sweat can be sampled and sensed non-invasively and continuously using electronic tattoos, bands, or patches. However, sweat as a diagnostic fluid presents numerous challenges as well, such as very small sample volumes and filtration (dilution) of larger-sized hydrophilic analytes. Currently the only major commercial application for sweat diagnostics is for infant cystic fibrosis testing based on sweat chloride concentrations.
| Biology and health sciences | Exocrine system | Biology |
77591 | https://en.wikipedia.org/wiki/Dianthus%20caryophyllus | Dianthus caryophyllus | Dianthus caryophyllus ( ), commonly known as carnation or clove pink, is a species of Dianthus native to the Mediterranean region. Its exact natural range is uncertain due to extensive cultivation over the last 2,000 years. Carnations are prized for their vibrant colors, delicate fringed petals, and fragrance.
The scent of carnations is often described as spicy, clove-like, or reminiscent of a combination of cinnamon and nutmeg, hence the common name "clove pink". This aroma has made carnations a popular choice for use in perfumes, potpourri, and scented products.
They have cultural significance and are associated with love, distinction, and motherly affection. With numerous cultivars and hybrids, carnations offer a wide variety of colors and forms, making them popular for gardens, floral arrangements, and scented products. Overall, carnations are enduring symbols of beauty and grace, treasured by flower enthusiasts and used to convey heartfelt emotions.
Taxonomy
Carnations were mentioned in Greek literature 2,000 years ago. The term dianthus was coined by Greek botanist Theophrastus, and is derived from the Ancient Greek words for divine ("dios") and flower ("anthos"). The name "carnation" is believed to come from the Latin , a "wreath, garland, chaplet, crown", as it was one of the flowers used in Greek and Roman ceremonial crowns, or possibly from the Latin (genitive ), "flesh", which refers to the natural colour of the flower, or in Christian iconography incarnatio, "incarnation", God made flesh in the form of Jesus.
Carl Linnaeus described the carnation in volume one of his Species Plantarum in 1753, giving it the name Dianthus caryophyllus.
Although originally applied to the species Dianthus caryophyllus, the name carnation is also often applied to some of the other species of Dianthus, and more particularly to garden hybrids between D. caryophyllus and other species in the genus.
Description
Dianthus caryophyllus is a herbaceous perennial plant growing up to 80 cm (3½ in) tall. (Correction needed) The leaves are glaucous greyish green to blue-green, slender, up to 15 cm (6 in) long. The flowers are produced singly or up to five together in a cyme; they are around 3-5 cm (1¼-2 in) diameter, and sweetly scented; the original natural flower color is bright pinkish-purple, but cultivars of other colors, including red, pink, yellow, white, and green have been developed. While sometimes dyed blue for cut bouquets, there are no known carnation cultivars that produce a true blue flower. The fragrant, hermaphrodite flowers have a radial symmetry. The four to six surrounding the calyx, egg-shaped, sting-pointed scales leaves are only ¼ as long as the calyx tube.
Distribution and habitat
The wild carnation is found in the Mediterranean countries of Portugal, Spain, Italy, Croatia, Albania, Greece and Turkey.
Cultivation and uses
Carnations require well-drained, neutral to slightly alkaline soil, and full sun. Numerous cultivars have been selected for garden planting. Typical examples include 'Gina Porto', 'Helen', 'Laced Romeo', and 'Red Rocket'. They are used for medical purposes, such as for upset stomach and fever. Their fragrance was historically used for vinegar, beer, wine, sauces and salads.
Crossbreeding D. caryophyllus with D. capitatus results in a hybrid that is resistant to bacterial wilt from Paraburkholderia caryophylli. However, the flower is less attractive and so more breeding and backcrossing is needed to improve the flower.
Carnation cultivars with no fragrance are often used by men as boutonnières or "buttonholes".
Diseases
Symbolism
Traditional meanings
For the most part, carnations express love, fascination, and distinction, though there are many variations dependent on color.
Along with the red rose, the red carnation can be used as a symbol of socialism and the labour movement, and historically has often been used in demonstrations on International Workers' Day (May Day).
In China, the carnation flower is the most frequently used flower in weddings.
In Portugal, bright red carnations were used when in 1974 the authoritarian Estado Novo regime was overthrown; therefore, this transition (brought about by a combination of a coup d'état with civil resistance) is known as the Carnation Revolution.
Light red carnations represent admiration, while dark red denote deep love and affection.
White carnations represent pure love and good luck, while striped (variegated) carnations symbolise regret that a love cannot be shared.
In the Netherlands, white carnations are associated with Prince Bernhard. He wore one during World War II and in a gesture of defiance some of the Dutch population took up this gesture. After the war the white carnation became a sign of the Prince, veterans and remembrance of the resistance.
Purple carnations indicate capriciousness. In France, it is a traditional funeral flower, given in condolence for the death of a loved one.
Carnation is the birth flower for those born in the month of January.
Since Ottoman times, red carnations and tulips are used in the interior wall paintings of mosques in Turkey. It is often said that while tulips represent God, carnations is the symbol for Muhammad. However these flower designs are not unique to mosques but also used in many other Ottoman traditional art forms.
In Azerbaijan, red carnations has turned into a symbol of mourning after the usage of the flower during Black January events of 1990, a violent crackdown on the civilian population of the country by USSR troops.
In the United States, former U.S. President William McKinley considered the red carnation to be his lucky flower and often wore one on his lapel. Following his assassination in 1901, the State of Ohio adopted the red carnation as its state flower to honor McKinley
The formal name for carnation, dianthus, comes from Greek for "heavenly flower", or the flower of Jove.
Holidays and events
Carnations are often worn on special occasions, especially Mother's Day and weddings. In 1907, Anna Jarvis chose a carnation as the emblem of Mother's Day because it was her mother's favourite flower. This tradition is now observed in the United States and Canada on the second Sunday in May. Ann Jarvis chose the white carnation because she wanted to represent the purity of a mother's love. This meaning has evolved over time, and now a red carnation may be worn if one's mother is alive, and a white one if she has died.
In Slovenia, red carnations are sometimes also given to women on Women's Day, the 8th of March, however, nowadays orchids or roses are often given as well.
In Korea, carnations express admiration, love and gratitude. Red and pink carnations are worn on Parents Day (Korea does not separate Mother's Day or Father's Day, but has Parents Day on 8 May). Sometimes, parents wear a corsage of carnation(s) on their left chest on Parents Day. Carnations are also worn on Teachers Day (15 May).
Red carnations are worn on May Day as a symbol of socialism and the labour movement in some countries, such as Austria, Italy, and successor countries of the former Yugoslavia. The red carnation is also the symbol of the Carnation Revolution in Portugal thanks to Celeste Caeiro.
Green carnations are for St. Patrick's Day and were famously worn by the Irish writer Oscar Wilde. The green carnation thence became a symbol of homosexuality in the early 20th century, especially through the book The Green Carnation and Noël Coward's song, "We All Wear a Green Carnation" in his operetta, Bitter Sweet.
In communist Czechoslovakia and in Poland in times of the People's Republic of Poland, carnations were traditionally given to women on the widely celebrated Women's Day, together with commodities that were difficult to obtain due to the countries' communist system, such as tights, towels, soap and coffee.
After the 1990 uprisings against Soviets in Azerbaijan in which 147 Azerbaijani civilians were killed, 800 people were injured and five people went missing, the carnation has become a symbol of the Black January tragedy associated with the carnations thrown into the puddles of blood shed in the streets of Azerbaijan subsequent to the massacre.
At the University of Oxford, carnations are traditionally worn to all examinations; white for the first exam, pink for exams in between, and red for the last exam. One story explaining this tradition relates that initially a white carnation was kept in a red inkpot between exams, so by the last exam it was fully red; the story is thought to originate in the late 1990s.
Carnations are the traditional first wedding anniversary flower.
Carnations are also known as the "Flower of God". The Greek name for Carnation is a fusion of "dios" and "anthos". Dios is used the described Zeus while Anthos means flower; thus the name "flower of God" is attached to it.
Symbols of territorial entities and organizations
The carnation is the national flower of Spain, Monaco, and Slovenia, and the provincial flower of the autonomous community of the Balearic Islands. The state flower of Ohio is a scarlet carnation, which was introduced to the state by Levi L. Lamborn. The choice was made to honor William McKinley, Ohio governor and U.S. president, who was assassinated in 1901, and regularly wore a scarlet carnation on his lapel.
Colours
Carnations do not naturally produce the pigment delphinidin, and thus a blue carnation cannot occur by natural selection or be created by traditional plant breeding. It shares this characteristic with other widely sold flowers like roses, lilies, tulips, chrysanthemums and gerberas.
Around 1996, a company, Florigene, used genetic engineering to extract certain genes from petunia and snapdragon flowers to produce a blue-mauve carnation, which was commercialized as Moondust. In 1998, a violet carnation called Moonshadow was commercialized. As of 2004, three additional blue-violet/purple varieties have been commercialized.
| Biology and health sciences | Caryophyllales | Plants |
77662 | https://en.wikipedia.org/wiki/Hagfish | Hagfish | Hagfish, of the class Myxini (also known as Hyperotreti) and order Myxiniformes , are eel-shaped jawless fish (occasionally called slime eels). Hagfish are the only known living animals that have a skull but no vertebral column, although they do have rudimentary vertebrae. Hagfish are marine predators and scavengers who can defend themselves against other larger predators by releasing copious amounts of slime from mucous glands in their skin.
Although their exact relationship to the only other living group of jawless fish, the lampreys, was long the subject of controversy, genetic evidence suggests that hagfish and lampreys are more closely related to each other than to jawed vertebrates, thus forming the clade Cyclostomi. The oldest-known stem group hagfish are known from the Late Carboniferous, around 310 million years ago, with modern representatives first being recorded in the mid-Cretaceous around 100 million years ago.
Physical characteristics
Body features
Hagfish are typically about in length. The largest-known species is Eptatretus goliath, with a specimen recorded at , while Myxine kuoi and Myxine pequenoi seem to reach no more than . Some have been seen as small as .
Hagfish have elongated, eel-like bodies, and paddle-like tails. The skin is naked and covers the body like a loosely fitting sock. They are generally a dull pink color and look quite worm-like. They have cartilaginous skulls (although the part surrounding the brain is composed primarily of a fibrous sheath) and tooth-like structures composed of keratin. Colors depend on the species, ranging from pink to blue-grey, and black or white spots may be present. Eyes are simple eyespots, not lensed eyes that can resolve images. Hagfish have no true fins and have six or eight barbels around the mouth and a single nostril. Instead of vertically articulating jaws like Gnathostomata (vertebrates with jaws), they have a pair of horizontally moving structures with tooth-like projections for pulling off food. The mouth of the hagfish has two pairs of horny, comb-shaped teeth on a cartilaginous plate that protracts and retracts. These teeth are used to grasp food and draw it toward the pharynx.
Its skin is attached to the body only along the center ridge of the back and at the slime glands, and is filled with close to a third of the body's blood volume, giving the impression of a blood-filled sack. It is assumed this is an adaptation to survive predator attacks. The Atlantic hagfish, representative of the subfamily Myxininae, and the Pacific hagfish, representative of the subfamily Eptatretinae, differ in that the latter has muscle fibers embedded in the skin. The resting position of the Pacific hagfish also tends to be coiled, while that of the Atlantic hagfish is stretched.
Slime
Hagfish can exude copious quantities of a milky and fibrous slime or mucus, from specialized slime glands. When released in seawater, the slime expands to 10,000 times its original size in 0.4 seconds. This slime that hagfish excrete has very thin fibers that make it more durable and retentive than the slime excreted by other animals. The fibers are made of proteins and also make the slime flexible. If they are caught by a predator, they can quickly release a large amount of slime to escape. If they remain captured, they can tie themselves in an overhand knot, and work their way from the head to the tail of the animal, scraping off the slime and freeing themselves from their captor. Rheological investigations showed that hagfish slime viscosity increases in elongational flow which favors gill clogging of suction feeding fish, while its viscosity decreases in shear which facilitates scraping off the slime by the travelling-knot.
Recently, the slime was reported to entrain water in its keratin-like intermediate filaments excreted by gland thread cells, creating a slow-to-dissipate, viscoelastic substance, rather than a simple gel. It has been shown to impair the function of a predator fish's gills. In this case, the hagfish's mucus would clog the predator's gills, disabling their ability to respire. The predator would release the hagfish to avoid suffocation. Because of the mucus, few marine predators target the hagfish. Other predators of hagfish are varieties of birds or mammals.
Free-swimming hagfish also slime when agitated, and later clear the mucus using the same travelling-knot behavior. The reported gill-clogging effect suggests that the travelling-knot behavior is useful or even necessary to restore the hagfish's own gill function after sliming.
Hagfish thread keratin (EsTKα and EsTKγ; and ), the protein that make up its slime filaments, is under investigation as an alternative to spider silk for use in applications such as body armor. These alpha-keratin proteins in hagfish slime transform from an α-helical structure to a stiffer β sheet structure when stretched. With combined draw-processing (stretching) and chemical crosslinking, recombinant slime keratin turns into a very strong fiber with an elastic modulus reaching 20 GPa.
When in 2017 a road accident on U.S. Highway 101 resulted in of hagfish being spilled, they emitted sufficient slime to cover the road and a nearby car.
Respiration
A hagfish generally respires by taking in water through its pharynx, past the velar chamber, and bringing the water through the internal gill pouches, which can vary in number from five to 16 pairs, depending on species. The gill pouches open individually, but in Myxine, the openings have coalesced, with canals running backwards from each opening under the skin, uniting to form a common aperture on the ventral side known as the branchial opening. The esophagus is also connected to the left branchial opening, which is therefore larger than the right one, through a pharyngocutaneous duct (esophageocutaneous duct), which has no respiratory tissue. This pharyngocutaneous duct is used to clear large particles from the pharynx, a function also partly taking place through the nasopharyngeal canal. In other species, the coalescence of the gill openings is less complete, and in Bdellostoma, each pouch opens separately to the outside, as in lampreys. The unidirectional water flow passing the gills is produced by rolling and unrolling velar folds located inside a chamber developed from the nasohypophyseal tract, and is operated by a complex set of muscles inserting into cartilages of the neurocranium, assisted by peristaltic contractions of the gill pouches and their ducts. Hagfish also have a well-developed dermal capillary network that supplies the skin with oxygen when the animal is buried in anoxic mud, as well as a high tolerance for both hypoxia and anoxia, with a well-developed anaerobic metabolism. Members of the group have spent 36 hours in water completely devoid of dissolved oxygen, and made a complete recovery. The skin has also been suggested to be capable of cutaneous respiration.
Nervous system
The origins of the vertebrate nervous system are of considerable interest to evolutionary biologists, and cyclostomes (hagfish and lampreys) are an important group for answering this question. The complexity of the hagfish brain has been an issue of debate since the late 19th century, with some morphologists suggesting that they do not possess a cerebellum, while others suggest that it is continuous with the midbrain. It is now considered that the hagfish neuroanatomy is similar to that of lampreys. A common feature of both cyclostomes is the absence of myelin in neurons. The brain of a hagfish has specific parts similar to the brains of other vertebrates. The dorsal and ventral muscles located towards the side of the hagfish body are connected to spinal nerves. The spinal nerves that connect to the muscles of the pharyngeal wall grow individually to reach them.
Eye
The hagfish eye lacks a lens, extraocular muscles, and the three motor cranial nerves (III, IV, and VI) found in more complex vertebrates, which is significant to the study of the evolution of more complex eyes. A parietal eye is also absent in extant hagfish. Hagfish eyespots, when present, can detect light, but as far as it is known, none can resolve detailed images. In Myxine and Neomyxine, the eyes are partly covered by the trunk musculature. Paleontological evidence suggests, however, that the hagfish eye is not plesiomorphic but rather degenerative, as fossils from the Carboniferous have revealed hagfish-like vertebrates with complex eyes. This would suggest that ancestrally Myxini possessed complex eyes.
Cardiac function, circulation, and fluid balance
Hagfish are known to have one of the lowest blood pressures among the vertebrates. One of the most primitive types of fluid balance found in animals is among these creatures; whenever a rise in extracellular fluid occurs, the blood pressure rises and this, in turn, is sensed by the kidney, which excretes excess fluid. They also have the highest blood volume to body mass of any chordate, with 17 ml of blood per 100 g of mass.
The hagfish circulatory system has been of considerable interest to evolutionary biologists and present day readers of physiology. Some observers first believed that the hagfish heart was not innervated (as the hearts of jawed vertebrates are), but further investigation revealed that the hagfish does have a true innervated heart. The hagfish circulatory system also includes multiple accessory pumps throughout the body, which are considered auxiliary "hearts".
Hagfish are the only known vertebrates with osmoregulation isosmotic to their external environment. Their renal function remains poorly described. There is a hypothesis that they excrete ions in bile salts.
Musculoskeletal system
Hagfish musculature differs from jawed vertebrates in that they have neither a horizontal septum nor a vertical septum, which in jawed vertebrates are junctions of connective tissue that separate the hypaxial musculature and epaxial musculature. They do, however, have true myomeres and myosepta like all vertebrates. The mechanics of their craniofacial muscles in feeding have been investigated, revealing advantages and disadvantages of their dental plate. In particular, hagfish muscles have increased force and gape size compared to similar-sized jawed vertebrates, but lack the speed amplification given by jawed vertebrates' muscles, suggesting that jaws are faster acting than hagfish dental plates.
The hagfish skeleton comprises the skull, the notochord, and the caudal fin rays. The first diagram of the hagfish endoskeleton was made by Frederick Cole in 1905. In Cole's monograph, he described sections of the skeleton that he termed "pseudo-cartilage", referring to its distinct properties compared to jawed chordates. The lingual apparatus of hagfish is composed of a cartilage base bearing two teeth-covered plates (dental plates) articulated with a series of large cartilage shafts. The nasal capsule is considerably expanded in hagfish, comprising a fibrous sheath lined with cartilage rings. In contrast to lampreys, the braincase is noncartilaginous. The role of their branchial arches is still highly speculative, as hagfish embryos undergo a caudal shift of the posterior pharyngeal pouches; thus, the branchial arches do not support gills. While parts of the hagfish skull are thought to be homologous with lampreys, they are thought to have very few elements homologous with jawed vertebrates.
Reproduction
Very little is known about hagfish reproduction. Obtaining embryos and observing reproductive behavior are difficult due to the deep-sea habitat of many hagfish species. In the wild, females outnumber males, with the exact sex-ratio differing depending on the species. E. burgeri, for example, has nearly a 1:1 ratio, while M. glutinosa females are significantly more common than males. Some species of hagfish are sexually undifferentiated before maturation, and possess gonadal tissue for both ovaries and testis. It has been suggested that females develop earlier than males, and that this may be the reason for unequal sex ratios. Hagfish testis are relatively small.
Depending on species, females lay from one to 30 tough, yolky eggs. These tend to aggregate due to having Velcro-like tufts at either end. It is unclear how hagfish go about laying eggs, although researchers have proposed three hypotheses based on observations of the low percentage of males and small testis. The hypotheses are that female hagfish lay eggs in small crevices in rock formations, the eggs are laid in burrow beneath the sand, and the slime produced by the hagfish is used to hold the eggs in a small area. It is worth noting that no direct evidence has been found to support any of these hypotheses. Hagfish do not have a larval stage, in contrast to lampreys.
Hagfish have a mesonephric kidney and are often neotenic of their pronephric kidney. The kidney(s) are drained via mesonephric/archinephric duct. Unlike many other vertebrates, this duct is separate from the reproductive tract, and the proximal tubule of the nephron is also connected with the coelom, providing lubrication. The single testicle or ovary has no transportation duct. Instead, the gametes are released into the coelom until they find their way to the posterior end of the caudal region, whereby they find an opening in the digestive system.
The hagfish embryo can develop for as long as 11 months before hatching, which is shorter in comparison to other jawless vertebrates. Not much was known about hagfish embryology until recently, when husbandry advances enabled considerable insight into the group's evolutionary development. New insights into the evolution of neural crest cells, support the consensus that all vertebrates share these cells, which might be regulated by a common subset of genes. Their genome has a large number of microchromosomes which are lost during the animal's development, leaving only the reproductive organs with a complete genome. Hagfish possess gonadotropins which secrete from pituitary glands to the gonads to stimulate development. This suggests that hagfish have an early version of the hypothalamic–pituitary–gonadal axis, a system which once thought to be exclusive to the Gnathostomes.Some species of hagfish reproduce seasonally, stimulated by hormones from their pituitary gland. E. burgeri is known to reproduce and migrate annually.
Feeding
While polychaete marine worms on or near the sea floor are a major food source, hagfish can feed upon and often even enter and eviscerate the bodies of dead and dying/injured sea creatures much larger than themselves. They are known to devour their prey from the inside. Hagfish have the ability to absorb dissolved organic matter across the skin and gill, which may be an adaptation to a scavenging lifestyle, allowing them to maximize sporadic opportunities for feeding. From an evolutionary perspective, hagfish represent a transitory state between the generalized nutrient absorption pathways of aquatic invertebrates and the more specialized digestive systems of aquatic vertebrates.
Like leeches, they have a sluggish metabolism and can survive months between feedings; their feeding behavior, however, appears quite vigorous. Analysis of the stomach content of several species has revealed a large variety of prey, including polychaetes, shrimp, hermit crabs, cephalopods, brittle stars, bony fishes, sharks, birds, and whale flesh.
In captivity, hagfish are observed to use the overhand-knot behavior in reverse (tail-to-head) to assist them in gaining mechanical advantage to pull out chunks of flesh from carrion fish or cetaceans, eventually making an opening to permit entry to the interior of the body cavity of larger carcasses. A healthy larger sea creature likely would be able to outfight or outswim this sort of assault.
This energetic opportunism on the part of the hagfish can be a great nuisance to fishermen, as they can devour or spoil entire deep drag-netted catches before they can be pulled to the surface. Since hagfish are typically found in large clusters on and near the bottom, a single trawler's catch could contain several dozen or even hundreds of hagfish as bycatch, and all the other struggling, captive sea life make easy prey for them.
The digestive tract of the hagfish is unique among chordates because the food in the gut is enclosed in a permeable membrane, analogous to the peritrophic matrix of insects. They are also able to absorb nutrients directly through their skin.
Hagfish have also been observed actively hunting the red bandfish, Cepola haastii, in its burrow, possibly using their slime to suffocate the fish before grasping it with their dental plates and dragging it from the burrow.
Classification
Originally, Myxine was included by Linnaeus (1758) in Vermes. The fossil hagfish Myxinikela siroka, from the Late Carboniferous of the United States, is the oldest-known member of the group. It is in some respects more similar to lampreys, but shows key autapomorphies of hagfish. In recent years, hagfish have become of special interest for genetic analysis investigating the relationships among chordates. Their classification as agnathans places hagfish as elementary vertebrates in between invertebrates and gnathostomes. However, discussion has long occurred in scientific literature about whether the hagfish were even invertebrate. Using fossil data, paleontologists posited that lampreys are more closely related to gnathostomes than hagfish. The term "Craniata" was used to refer to animals that had a developed skull, but were not considered true vertebrates. Molecular evidence in the early 1990s first began suggesting that lampreys and hagfish were more closely related to each other than to gnathostomes. The validity of the taxon "Craniata" was further examined by Delarbre et al. (2002) using mtDNA sequence data, concluding the Myxini are more closely related to the Hyperoartia than to the Gnathostomata—i.e., that modern jawless fishes form a clade called the Cyclostomata. The argument is that if the Cyclostomata are indeed monophyletic, Vertebrata would return to its old content (Gnathostomata + Cyclostomata) and the name Craniata, being superfluous, would become a junior synonym. Nowadays, molecular data are almost unanimously in consensus of cyclostome monophyly, with more recent work being directed at shared microRNAs between cyclostomes and gnathostomes. The current classification supported by molecular analyses (which show that lampreys and hagfishes are sister taxa), as well as the fact that hagfishes do, in fact, have rudimentary vertebrae, which places hagfishes in Cyclostomata.
Phylogeny
Hagfish are in the group Cyclostomata which includes jawless fish. The group Cyclostomata is characterized by two significant characteristics; keratinous tooth plates and movement of postotic myomeres to the orbitals. According to fossil record, hagfish and lampreys have been estimated to have diverged from one another during the Paleozoic period. An experiment used an estimation of synonymous and nonsynonymous substitutions for nucleotides and supplemented that data with pre-existing data into a clock that would calculate divergence times for the taxons Myxine and Eptatretus. This data found that the lineage diverged around 93–28 Mya; however, later studies have found even earlier divergence times within the group, with Myxine and Eptatretus diverging during the Triassic and the ancestor of Rubicundus diverging during other extant hagfishes during the Permian. Hagfish are excluded from the subphylum Gnathostomata because of morphological characteristics including the hagfish arched tongue. Hagfish embryos have characteristics of gnathostomes and may be plesiomorphic; however, these characteristics drastically change morphologically as the hagfish matures. The following hagfish and lamprey phylogeny is an adaptation based on the 2019 work of Miyashita et al.
Commercial use
As food
In most of the world, hagfish are not often eaten. But in Korea, the hagfish is a valued food, where it is generally skinned, coated in spicy sauce, and grilled over charcoal or stir-fried. It is especially popular in the southern port cities of the peninsula, such as Busan and coastal cities in South Gyeongsang Province.
Due to their value in Korean cuisine, most hagfish caught for food elsewhere in the world is fished with intent of being exported to South Korea. The inshore hagfish, found in the northwest Pacific, is eaten in Japan and South Korea. As hagfish slime binds vast amounts of liquid even at low temperatures, it was proposed as an energy-saving alternative for the production of tofu that does not require heating.
In textiles
The hagfish slime threads can be used as ultra-strong fiber for clothing. Douglas Fudge, of Chapman University, has conducted research in this area.
Skins
Hagfish skin, used in a variety of clothing accessories, is usually referred to as "eel skin". It produces a particularly durable leather, especially suitable for wallets and belts.
| Biology and health sciences | Agnatha | null |
77668 | https://en.wikipedia.org/wiki/Scar | Scar | A scar (or scar tissue) is an area of fibrous tissue that replaces normal skin after an injury. Scars result from the biological process of wound repair in the skin, as well as in other organs, and tissues of the body. Thus, scarring is a natural part of the healing process. With the exception of very minor lesions, every wound (e.g., after accident, disease, or surgery) results in some degree of scarring. An exception to this are animals with complete regeneration, which regrow tissue without scar formation.
Scar tissue is composed of the same protein (collagen) as the tissue that it replaces, but the fiber composition of the protein is different; instead of a random basketweave formation of the collagen fibers found in normal tissue, in fibrosis the collagen cross-links and forms a pronounced alignment in a single direction. This collagen scar tissue alignment is usually of inferior functional quality to the normal collagen randomised alignment. For example, scars in the skin are less resistant to ultraviolet radiation, and sweat glands and hair follicles do not grow back within scar tissues. A myocardial infarction, commonly known as a heart attack, causes scar formation in the heart muscle, which leads to loss of muscular power and possibly heart failure. However, there are some tissues (e.g. bone) that can heal without any structural or functional deterioration.
Types
All scarring is composed of the same collagen as the tissue it has replaced, but the composition of the scar tissue, compared to the normal tissue, is different. Scar tissue also lacks elasticity unlike normal tissue which distributes fiber elasticity. Scars differ in the amounts of collagen overexpressed. Labels have been applied to the differences in overexpression. Two of the most common types are hypertrophic and keloid scarring, both of which experience excessive stiff collagen bundled growth overextending the tissue, blocking off regeneration of tissues. Another form is atrophic scarring (sunken scarring), which also has an overexpression of collagen blocking regeneration. This scar type is sunken, because the collagen bundles do not overextend the tissue. Stretch marks (striae) are regarded as scars by some.
High melanin levels and either African or Asian ancestry may make adverse scarring more noticeable.
Hypertrophic
Hypertrophic scars occur when the body overproduces collagen, which causes the scar to be raised above the surrounding skin. Hypertrophic scars take the form of a red raised lump on the skin for lighter pigmented skin and the form of dark brown for darker pigmented skin. They usually occur within 4 to 8 weeks following wound infection or wound closure with excess tension and/or other traumatic skin injuries.
Keloid
Keloid scars are a more serious form of excessive scarring, because they can grow indefinitely into large, tumorous (although benign) neoplasms.
Hypertrophic scars are often distinguished from keloid scars by their lack of growth outside the original wound area, but this commonly taught distinction can lead to confusion.
Keloid scars can occur on anyone, but they are most common in dark-skinned people. They can be caused by surgery, cuts, accident, acne or, sometimes, body piercings. In some people, keloid scars form spontaneously. Although they can be a cosmetic problem, keloid scars are only inert masses of collagen and therefore completely harmless and not cancerous. However, they can be itchy or painful in some individuals. They tend to be most common on the shoulders and chest. Hypertrophic scars and keloids tend to be more common in wounds closed by secondary intention. Surgical removal of keloid is risky and may exacerbate the condition and worsening of the keloid.
Atrophic
An atrophic scar takes the form of a sunken recess in the skin, which has a pitted appearance. These are caused when underlying structures supporting the skin, such as fat or muscle, are lost. This type of scarring is often associated with acne, chickenpox, other diseases (especially Staphylococcus infection), surgery, certain insect and spider bites, or accidents. It can also be caused by a genetic connective tissue disorder, such as Ehlers–Danlos syndrome.
Stretch marks
Stretch marks (technically called striae) are also a form of scarring. These are caused when the skin is stretched rapidly (for instance during pregnancy, significant weight gain, or adolescent growth spurts), or when skin is put under tension during the healing process (usually near joints). This type of scar usually improves in appearance after a few years.
Elevated corticosteroid levels are implicated in striae development.
Umbilical
Humans and other placental mammals have an umbilical scar (commonly referred to as a navel) which starts to heal when the umbilical cord is cut after birth. Egg-laying animals have an umbilical scar which, depending on the species, may remain visible for life or disappear within a few days after birth.
Pathophysiology
A scar is the product of the body's repair mechanism after tissue injury. If a wound heals quickly within two weeks with new formation of skin, minimal collagen will be deposited and no scar will form. When the extracellular matrix senses elevated mechanical stress loading, tissue will scar, and scars can be limited by stress shielding wounds. Small full thickness wounds under 2mm reepithelize fast and heal scar free. Deep second-degree burns heal with scarring and hair loss. Sweat glands do not form in scar tissue, which impairs the regulation of body temperature. Elastic fibers are generally not detected in scar tissue younger than 3 months old. In scars, rete pegs are lost; through a lack of rete pegs, scars tend to shear easier than normal tissue.
The endometrium, the inner lining of the uterus, is the only adult tissue to undergo rapid cyclic shedding and regeneration without scarring, shedding and restoring roughly inside a 7-day window on a monthly basis. All other adult tissues, upon rapid shedding or injury, can scar.
Prolonged inflammation, as well as the fibroblast proliferation, can occur. Redness that often follows an injury to the skin is not a scar and is generally not permanent (see wound healing). The time it takes for this redness to dissipate may, however, range from a few days to, in some serious and rare cases, a few years.
Scars form differently based on the location of the injury on the body and the age of the person who was injured.
The worse the initial damage is, the worse the scar will generally be.
Skin scars occur when the dermis (the deep, thick layer of skin) is damaged. Most skin scars are flat and leave a trace of the original injury that caused them.
Wounds allowed to heal secondarily tend to scar worse than wounds from primary closure.
Collagen synthesis
An injury does not become a scar until the wound has completely healed; this can take many months, or years in the worst pathological cases, such as keloids. To begin to patch the damage, a clot is created; this clot is the beginning process that results in a provisional matrix. In the process, the first layer is a provisional matrix and is not a scar. Over time, the wounded body tissue overexpresses collagen inside the provisional matrix to create a collagen matrix. This collagen overexpression continues and crosslinks the fiber arrangement inside the collagen matrix, making the collagen dense. This densely packed collagen, morphing into an inelastic whitish collagen scar wall, blocks off cell communication and regeneration; as a result, the new tissue generated will have a different texture and quality than the surrounding unwounded tissue. This prolonged collagen-producing process results in a fortuna scar.
Fibroblasts
The scarring is created by fibroblast proliferation, a process that begins with a reaction to the clot. To mend the damage, fibroblasts slowly form the collagen scar. The fibroblast proliferation is circular and cyclically, the fibroblast proliferation lays down thick, whitish collagen inside the provisional and collagen matrix, resulting in the abundant production of packed collagen on the fibers giving scars their uneven texture. Over time, the fibroblasts continue to crawl around the matrix, adjusting more fibers and, in the process, the scarring settles and becomes stiff. This fibroblast proliferation also contracts the tissue. In unwounded tissue, these fibers are not overexpressed with thick collagen and do not contract.
EPF and ENF fibroblasts have been genetically traced with the Engrailed-1 genetic marker. EPFs are the primary contributors to all fibrotic outcomes after wounding. ENFs do not contribute to fibrotic outcomes.
Myofibroblast
Mammalian wounds that involve the dermis of the skin heal by repair, not regeneration (except in 1st trimester inter-uterine wounds and in the regeneration of deer antlers). Full-thickness wounds heal by a combination of wound contracture and edge re-epitheliasation. Partial thickness wounds heal by edge re-epithelialisation and epidermal migration from adnexal structures (hair follicles, sweat glands and sebaceous glands). The site of keratinocyte stem cells remains unknown but stem cells are likely to reside in the basal layer of the epidermis and below the bulge area of hair follicles.
The fibroblast involved in scarring and contraction is the myofibroblast, which is a specialized contractile fibroblast. These cells express α-smooth muscle actin (α-SMA).
The myofibroblasts are absent in the first trimester in the embryonic stage where damage heals scar-free; in small incisional or excision wounds less than 2 mm that also heal without scarring; and in adult unwounded tissues where the fibroblast in itself is arrested; however, the myofibroblast is found in massive numbers in adult wound healing which heals with a scar.
The myofibroblasts make up a high proportion of the fibroblasts proliferating in the postembryonic wound at the onset of healing. In the rat model, for instance, myofibroblasts can constitute up to 70% of the fibroblasts, and is responsible for fibrosis on tissue. Generally, the myofibroblasts disappear from the wound within 30 days, but can remain in pathological cases in hypertrophy, such as keloids. Myofibroblasts have plasticity and in mice can be transformed into fat cells, instead of scar tissue, via the regeneration of hair follicles.
Mechanical stress
Wounds under 2mm generally do not scar but larger wounds generally do scar. In 2011 it was found that mechanical stress can stimulate scarring and that stress shielding can reduce scarring in wounds. In 2021 it was found that using chemicals to manipulate fibroblasts to not sense mechanical stress brought scar-free healing. The scar-free healing also occurred when mechanical stress was placed onto a wound.
Treatment
Early and effective treatment of acne scarring can prevent severe acne and the scarring that often follows. In 2004, no prescription drugs for the treatment or prevention of scars were available.
Chemical peels
Chemical peels are chemicals which destroy the epidermis in a controlled manner, leading to exfoliation and the alleviation of certain skin conditions, including superficial acne scars. Various chemicals can be used depending upon the depth of the peel, and caution should be used, particularly for dark-skinned individuals and those individuals susceptible to keloid formation or with active infections.
Filler injections
Filler injections of collagen can be used to raise atrophic scars to the level of surrounding skin. Risks vary based upon the filler used, and can include further disfigurement and allergic reaction.
Laser treatment
Nonablative lasers, such as the 585 nm pulsed dye laser, 1064 nm and 1320 nm Nd:YAG, or the 1540 nm Er:Glass are used as laser therapy for hypertrophic scars and keloids. There is tentative evidence for burn scars that they improve the appearance.
Ablative lasers such as the carbon dioxide laser (CO) or Er:YAG offer the best results for atrophic and acne scars. Like dermabrasion, ablative lasers work by removing the epidermis. Healing times for ablative therapy are much longer and the risk profile is greater compared to nonablative therapy; however, nonablative therapy offers only minor improvements in cosmetic appearance of atrophic and acne scars.
Radiotherapy
Low-dose, superficial radiotherapy is sometimes used to prevent recurrence of severe keloid and hypertrophic scarring. It is thought to be effective despite a lack of clinical trials, but only used in extreme cases due to the perceived risk of long-term side effects.
Dressings and topical silicone
Silicone scar treatments are commonly used in preventing scar formation and improving existing scar appearance. A meta-study by the Cochrane collaboration found weak evidence that silicone gel sheeting helps prevent scarring. However, the studies examining it were of poor quality and susceptible to bias.
Pressure dressings are commonly used in managing burn and hypertrophic scars, although supporting evidence is lacking. Care providers commonly report improvements, however, and pressure therapy has been effective in treating ear keloids. The general acceptance of the treatment as effective may prevent it from being further studied in clinical trials.
Verapamil-containing silicone gel
Verapamil, a type of calcium channel blocker, is considered a candidate drug for the treatment of hypertrophic scars. A study conducted by the Catholic University of Korea concluded that verapamil-releasing silicone gel is effective and is a superior alternative to the conventional silicone gel where decreased median SEI, fibroblast count, and collagen density in all verapamil-added treatment groups were observed. Gross morphologic features suggested that the combination of verapamil and silicone improves the overall quality of hypertrophic scars by reducing scar height and redness. This was verified with quantifiable histomorphometric parameters; however, oral verapamil is not a good choice because of its effect of lowering blood pressure. Intralesional injection of verapamil is also suboptimal because of the required frequency for injections. Topical silicone gel combined with verapamil does not lead to systemic hypotension, is convenient to apply, and shows enhanced results.
Steroids
A long-term course of corticosteroid injections into the scar may help flatten and soften the appearance of keloid or hypertrophic scars.
Topical steroids are ineffective. However, clobetasol propionate can be used as an alternative treatment for keloid scars.
Topical steroid applied immediately after fractionated laser treatment is however very effective (and more efficacious than laser treatment alone) and has shown benefit in numerous clinical studies.
Surgery
Scar revision is a process of cutting the scar tissue out. After the excision, the new wound is usually closed up to heal by primary intention, instead of secondary intention. Deeper cuts need a multilayered closure to heal optimally, otherwise depressed or dented scars can result.
Surgical excision of hypertrophic or keloid scars is often associated to other methods, such as pressotherapy or silicone gel sheeting. Lone excision of keloid scars, however, shows a recurrence rate close to 45%. A clinical study is currently ongoing to assess the benefits of a treatment combining surgery and laser-assisted healing in hypertrophic or keloid scars.
Subcision is a process used to treat deep rolling scars left behind by acne or other skin diseases. It is also used to lessen the appearance of severe glabella lines, though its effectiveness in this application is debatable. Essentially the process involves separating the skin tissue in the affected area from the deeper scar tissue. This allows the blood to pool under the affected area, eventually causing the deep rolling scar to level off with the rest of the skin area. Once the skin has leveled, treatments such as laser resurfacing, microdermabrasion or chemical peels can be used to smooth out the scarred tissue.
Vitamins
Research shows the use of vitamin E and onion extract (sold as Mederma) as treatments for scars is ineffective. Vitamin E causes contact dermatitis in up to 33% of users and in some cases it may worsen scar appearance and could cause minor skin irritations, but Vitamin C and some of its esters fade the dark pigment associated with some scars.
Other
Cosmetics; Medical makeup can temporarily conceal scars. This is most commonly used for facial scars.
Dermabrasion involves the removal of the surface of the skin with special equipment, and usually involves a local anaesthetic.
A 2012 literature review found weak evidence that massage was efficacious in scar management. Any beneficial effect appeared to be greater in wounds created by surgical incision than for traumatic or burn wounds. A 2022 scoping review covering twenty-five studies of 1515 participants reported that all studies reviewed reported favorable outcomes for scar massage, but that "while there may be benefits to scar massage in reducing pain, increasing movement and improving scar characteristics", there was a lack of "consistent research methods, intervention protocols and outcome measures".
Microneedling
Society and culture
Intentional scarring
The permanence of scarring has led to its intentional use as a form of body art within some cultures and subcultures. These forms of ritual and non-ritual scarring practices can be found in many groups and cultures around the world.
Etymology
First attested in English in the late 14th century, the word scar derives from a conflation of Old French escharre, from Late Latin eschara, which is the Latinisation of the Greek ἐσχάρα (eskhara), meaning "hearth, fireplace", but in medicine "scab, eschar on a wound caused by burning or otherwise", and Middle English skar ("cut, crack, incision"), which is from Old Norse skarð ("notch, gap"). The conflation helped to form the English meaning. Compare the place name Scarborough for evolution of skarð to scar.
Research
Research, before 2009, focused on scar improvements with research into molecular mechanisms. Treatments involving molecular mechanisms including avotermin, ribosomal s6 kinase (RSK), and osteopontin were investigated at the time. After successful phase I/II trials, human recombinant TGF-β3 (avotermin, planned trade name Juvista) failed in Phase III trials.
In 2011, the scientific literature highlighted stress shielding a fresh wound through the wound healing process, brings significant scar improvement and smaller scars.
By 2016, skin had been regenerated in vivo and in vitro. and scar-free healing had been operationalized and induced by four main regeneration techniques: by instrument, by materials, by drugs, and by in vitro 3-D printing. In 2018, a silk-derived sericin hydrogel dressing was undergoing research, the material was shown to prevent scar formation. By 2021, more people were paying attention to the possibility of scar revision and new technologies.
In 2021, researchers found that, verteporfin, an FDA-approved drug for eye disease, could enable scar-free healing in mice. According to the study, the drug works by blocking mechanical stress signals in fibroblast cells.
| Biology and health sciences | Injury: General | Health |
77688 | https://en.wikipedia.org/wiki/Attosecond | Attosecond | An attosecond (abbreviated as as) is a unit of time in the International System of Units (SI) equal to 10−18 or 1⁄1 000 000 000 000 000 000 (one quintillionth) of a second.
An attosecond is to a second, as a second is to approximately 31.69 billion years.
The attosecond is a tiny unit but it has various potential applications: it can observe oscillating molecules, the chemical bonds formed by atoms in chemical reactions, and other extremely tiny and extremely fast things.
Common measurements
0.247 attoseconds: travel time of a photon across "the average bond length of molecular hydrogen"
24.189... attoseconds: the atomic unit of time
43 attoseconds: the shortest pulses of laser light yet created
53 attoseconds: the shortest electron laser pulse ever created
53 attoseconds: the second-shortest pulses of laser light created
82 attoseconds (approximately): half-life of beryllium-8, maximum time available for the triple-alpha process for the synthesis of carbon and heavier elements in stars
84 attoseconds: the approximate half-life of a neutral pion
100 attoseconds: fastest-ever view of molecular motion
320 attoseconds: the estimated time it takes electrons to transfer between atoms
Historical development
In 2001, Ferenc Krausz and his team at the Technical University of Vienna fired an ultrashort wavelength (7 femtoseconds) red laser pulse into a stream of neon atoms, where the stripped electrons were carried by the pulse and almost immediately re-eject into the neon nucleus.
While capturing the attosecond pulse, the physicists also demonstrated its utility. They aimed attosecond and longer-wavelength red pulses at a type of krypton atom simultaneously: first, the electrons were knocked off; then, the red light pulse hit the electrons; finally, the energy was tested. Judging from the difference in the timing of these two pulses, the scientists obtained a very precise measurement of how long it took the electron to decay (how many attoseconds). Never before have scientists used such a short time scale to study the energy of electrons.
Applications
Need for more precise units
The crystal lattice vibrates and molecules rotate on a scale of picoseconds. The creation and breaking of chemical bonds and molecular vibration happen in femtoseconds. Observing the motion of electrons happens on the attosecond scale.
The number of electrons in an atom and their configuration define an element. Because attosecond pulses are faster than the motion of electrons in atoms and molecules, attosecond provides a new tool for controlling and measuring quantum states of matter. These pulses have been used to explore the detailed physics of atoms and molecules and have potential applications in fields ranging from electronics to medicine.
Directly observing the wave oscillations of light
Using a method called attosecond streaking, people can see the electrical components of EM waves. Scientists start with a gas of neon atoms and ionize them with a single ultrashort burst of UV radiation measured in attoseconds. The electric field of the infrared can then strongly influence the motion of the electrons. The electrons will be forced up and down as the field oscillates. Depending on when the electron is released, this process will emit different final energies. The final measurement of the electron's energy, as a function of the relative delay between the two pulses, clearly shows the traces of the electric field of the attosecond pulse.
Short pulses of light
The 2023 Nobel Prize in Physics was awarded to Pierre Agostini, Ferenc Krausz, and Anne L'Huillier for demonstrating a way to create "almost unimaginably" short pulses of light, measured in attoseconds. These pulses can be used to capture and study rapid processes inside atoms, such as the behavior of electrons.
| Physical sciences | Time | Basics and measurement |
77689 | https://en.wikipedia.org/wiki/Femtosecond | Femtosecond | A femtosecond is a unit of time in the International System of Units (SI) equal to 10 or of a second; that is, one quadrillionth, or one millionth of one billionth, of a second.
A femtosecond is to a second, as a second is to approximately 31.69 million years.
For context, a ray of light travels approximately 0.3 μm (micrometers) in 1 femtosecond, a distance comparable to the diameter of a virus. The first to make femtosecond measurements was the Egyptian Nobel Laureate Ahmed Zewail, for which he was awarded the Nobel Prize in Chemistry in 1999. Professor Zewail used lasers to measure the movement of particles at the femtosecond scale, thereby allowing chemical reactions to be observed for the first time.
The word femtosecond is formed by the SI prefix femto and the SI unit second. Its symbol is fs.
A femtosecond is equal to 1000 attoseconds, or 1/1000 picosecond. Because the next higher SI unit is 1000 times larger, times of 10−14 and 10−13 seconds are typically expressed as tens or hundreds of femtoseconds.
Typical time steps for molecular dynamics simulations are on the order of 1 fs.
The periods of the waves of visible light have a duration of about 2 femtoseconds. The precise duration depends on the energy of the photons, which determines their color. (See wave–particle duality.) This time can be calculated by dividing the wavelength of the light by the speed of light (approximately 3×108 m/s) to determine the time required for light to travel that distance.
Examples
46 fs – the swiftest chemical reaction known (radiolysis of water leads to the formation of a H2O+ ion, which rapidly reacts to become hydronium (H3O+) and a short lived hydroxyl radical (•OH))
200 fs – the average chemical reaction, such as the reaction of pigments in an eye to light
300 fs – the duration of a vibration of the atoms in an iodine molecule
| Physical sciences | Time | Basics and measurement |
77700 | https://en.wikipedia.org/wiki/Picosecond | Picosecond | A picosecond (abbreviated as ps) is a unit of time in the International System of Units (SI) equal to 10−12 or (one trillionth) of a second. That is one trillionth, or one millionth of one millionth of a second, or 0.000 000 000 001 seconds.
A picosecond is to one second, as one second is to approximately 31,688.76 years.
Multiple technical approaches achieve imaging within single-digit picoseconds: for example, the streak camera or intensified CCD (ICCD) cameras are able to picture the motion of light.
One picosecond is equal to 1000 femtoseconds, or 1/1000 nanoseconds. Because the next SI unit is 1000 times larger, measurements of 10−11 and 10−10 second are typically expressed as tens or hundreds of picoseconds. Some notable measurements in this range include:
1.0 picoseconds (1.0 ps) – cycle time for electromagnetic frequency 1 terahertz (THz) (1 x 1012 hertz), an inverse unit. This corresponds to a wavelength of 0.3 mm, as can be calculated by multiplying 1 ps by the speed of light (approximately 3 x 108 m/s) to determine the distance traveled. 1 THz is in the far infrared.
1 picosecond – time taken by light in vacuum to travel approximately 0.30 mm
1 picosecond – half-life of a bottom quark
~1 picosecond – lifetime of a single (hydronium) ion in water at 20 °C
picoseconds to nanoseconds – phenomena observable by dielectric spectroscopy
1.2 picoseconds – switching time of the world's fastest transistor (845 GHz, as of 2006)
1.7 picoseconds – rotational correlation time of water
3.3 picoseconds (approximately) – time taken for light to travel 1 millimeter
10 picoseconds after the Big Bang – electromagnetism separates from the other fundamental forces
34 picoseconds – signal rise time (20% to 80%) of a SFP+ transmitter for 10 Gigabit Ethernet.
10 to 150 picoseconds – rotational correlation times of a molecule (184 g/mol) from hot to frozen water
100 picoseconds – Unit Interval of a 10 Gbit/s serial communication link, such as USB 3.1.
108.7827757 picoseconds – transition time between the two hyperfine levels of the ground state of the caesium-133 atom at absolute zero
330 picoseconds (approximately) – the time it takes a common 3.0 GHz computer CPU to complete a processing cycle
| Physical sciences | Time | Basics and measurement |
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