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77895 | https://en.wikipedia.org/wiki/Indri | Indri | The indri (; Indri indri), also called the babakoto, is one of the largest living lemurs, with a head-body length of about and a weight of between . It has a black and white coat and maintains an upright posture when climbing or clinging. It is monogamous and lives in small family groups, moving through the canopy, and is herbivorous, feeding mainly on leaves but also seeds, fruits, and flowers. The groups are quite vocal, communicating with other groups by singing, roaring and other vocalisations. Besides humans, it is the only mammal found that can use rhythm.
It is a diurnal tree-dweller related to the sifakas and, like all lemurs, it is native to Madagascar. It is revered by the Malagasy people and plays an important part in their myths and legends with various stories in existence accounting for its origin. The main threats faced by the indri are habitat destruction and fragmentation due to slash and burn agriculture, fuelwood gathering, and logging. It is also hunted despite taboos against this. The International Union for Conservation of Nature has rated its conservation status as "critically endangered".
Etymology
The name "indri" most likely comes from a native Malagasy name for the animal, endrina. An oft-repeated, but incorrect story is that the name comes from indry , meaning "there" or "there it is".
French naturalist Pierre Sonnerat, who first described the animal, supposedly heard a Malagasy point out the animal and took the word to be its name. It has been suggested that he may have heard the local name endrina which is used. Another Malagasy name for the animal is . Babakoto is most commonly translated as "ancestor" or "father", but several translations are possible. "Koto" is a Malagasy word for "little boy" and "Rakoto" is a common name, with "Koto" as its diminutive. As "baba" is a term for "father," the word "babakoto" may be translated as "father of a little boy" or "father of Rakoto." The father-son dynamic of many of the babakoto origin myths helps to explain the Malagasy name.
Physical characteristics
Along with the diademed sifaka, the indri is the largest lemur still in existence; both have average weights of about 6.5 kg. It can weigh up to to and perhaps up to . It has a head-body length of and can reach nearly with legs fully extended.
The indri is a vertical clinger and leaper and thus holds its body upright when traveling through trees or resting in branches. It has long, muscular legs which it uses to propel itself from trunk to trunk. Its large greenish eyes and black face are framed by round, fuzzy ears. Unlike any other living lemur, the indri has only a rudimentary tail. The silky fur is mostly black with white patches along the limbs, neck, crown, and lower back. Different populations of the species show wide variations in color, with some northern populations consisting of mostly or entirely black individuals. The face is bare with pale black skin, and it is sometimes fringed with white fur.
Due to these color variations, Colin Groves listed two subspecies of the indri in 2005: The dark Indri indri indri from the northern part of its range and the relatively pale Indri indri variegatus from the southern part. Later editions of Lemurs of Madagascar by Russell Mittermeier et al. do not recognize this classification, and recent genetic and morphological work suggests the variation in the indri is clinal.
Behavior
The indri practices long-term monogamy, seeking a new partner only after the death of a mate. It lives in small groups consisting of the mated male and female and their maturing offspring. In the more fragmented forests of their range, the indri may live in larger groups with several generations. Habitat fragmentation limits the mobility and capacity of these large groups to break into smaller units.
Like many other species of lemur, indri live in a female dominant society. The dominant female often will displace males to lower branches and poorer feeding grounds, and is typically the one to lead the group during travel.
It is common for groups to move 300–700 m daily, with most distance travelled midsummer in search of fruit. Indris sleep in trees about 10–30 m above ground and typically sleep alone or in pairs. It is common for young female indris, occasionally adult females, to silently play wrestle anywhere from a few seconds up to 15 minutes. Members of a single group will urinate and defecate jointly at one of their many selected areas of defecation in their territory.
Reproduction
Indris reach sexual maturity between the ages of 7 and 9. Females bear offspring every two to three years, with a gestation period around 120–150 days. The single infant is usually born in May or June. The mother is the primary caregiver, though the father assists, remaining with his mate and offspring. Infants are born mostly or completely black and begin to show white coloration (if any) between four and six months of age. The infant clings to its mother's belly until it is four or five months old, at which time it is ready to move onto her back. The indri begins to demonstrate independence at eight months, but it will not be fully independent from its mother until it is at least two years old.
Communication
The indri makes loud, distinctive songs, which can last from 45 seconds to more than 3 minutes. Song duration and structure varies among and even within groups, but most songs have the following three-phase pattern.
Usually, a roaring sequence lasting for several seconds will precede the more characteristic vocalizations. All members of the group except the very young participate in this roar, but the song proper is dominated by the adult pair. They follow the roar with a long note sequence, characterized by notes of up to five seconds in duration. After this is a descending phrase sequence. The wails begin on a high note and become progressively lower-pitched. It is common for two or more indri to coordinate the timing of their descending notes to form a duet.
Different indri groups typically sing sequentially, responding to one another. As well as solidifying contacts between groups, the songs may communicate territorial defense and boundaries, environmental conditions, reproductive potential of the group members, and warning signals. The indri may sing after disturbances such as thunder, airplanes, bird calls, and other lemur calls. A group will sing almost every day, up to seven times daily. The peak singing hours are between 7 and 11 am. Daily frequency of song is highest during the indri's breeding season from December to March.
Several other indri vocalizations have been identified. The "roar" is also used as a warning signal for aerial predators such as hawks. The indri emit a "hoot" or "honk" to warn of terrestrial predators such as the fossa. Other vocal categories include the "grunt", "kiss", "wheeze", and "hum". The purpose of these is not well understood.
Before singing, the indri move to the tree tops, which allows them to be heard up to 4 km away.
Diet and feeding
The indri is herbivorous and primarily folivorous. It prefers young, tender leaves, but will also eat seeds, fruits, and flowers. Female indri seem to have greater preference for immature leaves than males do and spend more time foraging among them. A wide variety of plant species are consumed, with members of the laurel family featuring prominently in the diet. The indri consumes little nontree vegetation.
To feed, the indri plucks off a leaf or other plant part with its teeth. It uses its hands to pull tree branches closer to its mouth.
Reproductively mature females have priority access to food sources, therefore they forage higher in the trees than males.
Distribution
This lemur inhabits the lowland and montane forests along the eastern coast of Madagascar, from the Réserve Spéciale d’Anjanaharibe-Sud in the north to the Mangoro River in the south. They are absent from the Masoala Peninsula and the Marojejy National Park, even though both regions are connected to forests where indri do occur less than 40 km away.
Relationship with humans
Mythology
Across Madagascar, the indri is revered and protected by fady (taboos). Countless variations are given on the legend of the indri's origins, but they all treat it as a sacred animal, not to be hunted or harmed.
A legend tells of a man who went hunting in the forest and did not return. His absence worried his son, who went out looking for him. When the son also disappeared, the rest of the villagers ventured into the forest seeking the two, but discovered only two large lemurs sitting in the trees: the first indri. The boy and his father had transformed. In some versions, only the son transforms, and the wailing of the babakoto is analogous to the father's wailing for his lost son.
Another human-like characteristic of the indri is its behavior in the sun. Like its sifaka relatives, the indri frequently engages in what has been described as sun-bathing or sun-worshipping. As the sun rises each morning, it will sit and face it from a tree branch with its legs crossed, back straight, hands low with palms facing out or resting on its knees, and eyes half-closed. Biologists are hesitant to call this behavior sun worship, as the term may be overly anthropomorphic. However, many Malagasy people do believe that the indri worships the sun.
Conservation
The first film of indri was obtained using tape lures, on an expedition forming the basis of David Attenborough's 1961 BBC series Zoo Quest to Madagascar.
The indri is a critically endangered species. While population estimates are uncertain (1,000 – 10 000 individuals), the population appears to be rapidly shrinking and may diminish by 80% over the next three generations (~36 years). The primary threats to its existence are habitat destruction and fragmentation due to slash and burn agriculture, fuelwood gathering, and logging. This kind of destruction occurs even in protected areas.
The indri is also widely hunted, despite the many origin myths and traditional taboos (fady) that hold it sacred. Cultural erosion and immigration are partly to blame for the breakdown of traditional beliefs. In some cases, Malagasy people who resent the protective fady find ways to circumvent them. People whose fady forbid them from eating the indri may still hunt the lemurs and sell their flesh, and those forbidden to kill the indri may still purchase and consume them. Indri meat is prized as a delicacy in some regions.
Only one indri has lived over a year in captivity and none have bred successfully while captive.
| Biology and health sciences | Strepsirrhini | Animals |
78000 | https://en.wikipedia.org/wiki/Tablet%20%28pharmacy%29 | Tablet (pharmacy) | A tablet (also known as a pill) is a pharmaceutical oral dosage form (oral solid dosage, or OSD) or solid unit dosage form. Tablets may be defined as the solid unit dosage form of medication with suitable excipients. It comprises a mixture of active substances and excipients, usually in powder form, that are pressed or compacted into a solid dose. The main advantages of tablets are that they ensure a consistent dose of medicine that is easy to consume.
Tablets are prepared either by moulding or by compression. The excipients can include diluents, binders or granulating agents, glidants (flow aids) and lubricants to ensure efficient tabletting; disintegrants to promote tablet break-up in the digestive tract; sweeteners or flavours to enhance taste; and pigments to make the tablets visually attractive or aid in visual identification of an unknown tablet. A polymer coating is often applied to make the tablet smoother and easier to swallow, to control the release rate of the active ingredient, to make it more resistant to the environment (extending its shelf life), or to enhance the tablet's appearance.
Medicinal tablets were originally made in the shape of a disk of whatever colour their components determined, but are now made in many shapes and colours to help distinguish different medicines. Tablets are often imprinted with symbols, letters, and numbers, which allow them to be identified, or a groove to allow splitting by hand. Sizes of tablets to be swallowed range from a few millimetres to about a centimetre.
The compressed tablet is the most commonly seen dosage form in use today. About two-thirds of all prescriptions are dispensed as solid dosage forms, and half of these are compressed tablets. A tablet can be formulated to deliver an accurate dosage to a specific site in the body; it is usually taken orally, but can be administered sublingually, buccally, rectally or intravaginally. The tablet is just one of the many forms that an oral drug can take such as syrups, elixirs, suspensions, and emulsions.
History
Pills are thought to date back to around 1500 BC. Earlier medical recipes, such as those from 4000 BC, were for liquid preparations rather than solids. The first references to pills were found on papyruses in ancient Egypt and contained bread dough, honey, or grease. Medicinal ingredients, such as plant powders or spices, were mixed in and formed by hand to make little balls, or pills.
In ancient Greece, such medicines were known as katapotia ("something to be swallowed"), and the Roman scholar Pliny, who lived from 23 to 79 AD, first gave a name to what we now call pills, calling them pilula.
Pills have always been difficult to swallow, and efforts have been made to make them go down easier. In mediaeval times, people coated pills with slippery plant substances. Another approach, used as recently as the 19th century, was to gild them in gold and silver, although this often meant that they would pass through the digestive tract with no effect. In the 1800s, sugar coating and gelatin coating were invented, as were gelatin capsules.
In 1843, the British painter and inventor William Brockedon was granted a patent for a machine capable of "Shaping Pills, Lozenges, and Black Lead by Pressure in Dies". The device was capable of compressing powder into a tablet without the use of an adhesive.
Types
Pills
A pill was originally defined as a small, round, solid pharmaceutical oral dosage form of medication. The word's etymology reflects the historical concepts of grinding the ingredients with a mortar and pestle and rolling the resultant paste or dough into lumps to be dried. Today, in its strict sense, the word pill still refers specifically to tablets (including caplets) rather than capsules (which were invented much later), but because a simple hypernym is needed to intuitively cover all such oral dosage forms, the broad sense of the word pill is also widely used and includes both tablets and capsules — colloquially, any solid oral form of medication falls into the "pill" category (see pill § Usage notes).
An early example of a pill comes from ancient Rome. They were made of zinc carbonates, hydrozincite and smithsonite. The pills were used for sore eyes and were found aboard a Roman ship that wrecked in 140 BC. However, these tablets were meant to be pressed on the eyes, not swallowed.
Defects/imperfections arising during tablet manufacturing
Formulation related: sticking, picking, binding
Processing: capping, lamination, cracking, chipping
Machine: double impression
Caplets
A caplet is a smooth, coated, oval-shaped medicinal tablet in the general shape of a capsule. Many caplets have an indentation running down the middle, so they may be split in half more easily. Consumers have viewed capsules as the most effective way to take medication ever since they first appeared. For this reason, producers of drugs such as OTC analgesics wanting to emphasize the strength of their product developed the "caplet", a portmanteau of capsule-shaped tablet, in order to tie this positive association to more efficiently produced tablet pills as well as being an easier-to-swallow shape than the usual disk-shaped tablet.
Orally disintegrating tablets (ODT)
An orally disintegrating tablet or orodispersible tablet (ODT), is a drug dosage form available for a limited range of over-the-counter (OTC) and prescription medications.
Tablet formulations
In the tablet-pressing process, it is important that all ingredients be fairly dry, powdered or granular, somewhat uniform in particle size, and freely flowing. Mixed particle sized powders may separate during the manufacturing process due to differing particle densities. This can result in tablets with non-uniform concentrations of drug or active pharmaceutical ingredient (API), resulting in uneven dosage between tablets, but granulation should prevent this.
Some APIs may be compressed into tablets as pure substances, but this is rarely the case; most formulations include pharmacologically inactive ingredients (excipients):
Usually a binder is added to help hold the tablet together and give it mechanical strength. A wide variety of binders may be used: some common ones are lactose, dibasic calcium phosphate, sucrose, corn (maize) starch, microcrystalline cellulose, povidone polyvinylpyrrolidone and modified cellulose (for example, hydroxypropyl methylcellulose and hydroxyethylcellulose).
Often, an ingredient is also needed to act as a disintegrator to aid tablet dispersion once swallowed, releasing the API for absorption. Some binders, such as starch and cellulose, are also excellent disintegrators.
Tablet properties
Tablets can be made in virtually any shape, although the requirements of patients and tableting machines mean that most are round, oval, or capsule-shaped. More unusual shapes have been manufactured, but patients find these harder to swallow, and they are more vulnerable to chipping or manufacturing problems.
Tablet diameter and shape are determined by the machine tooling used to produce them; a die plus an upper and a lower punch are required. This is called a station of tooling. The amount of tablet material and the placement of the punches in relation to one another during compression determine the thickness. Once this is done, we can measure the corresponding pressure applied during compression. The shorter the distance between the punches, the greater the pressure applied during compression, and sometimes the harder the tablet. Tablets need to be hard enough that they do not break up in the bottle, yet friable enough that they disintegrate in the gastric tract.
Tablets need to be strong enough to resist the stresses of packaging, shipping, and handling by the pharmacist and patient. The mechanical strength of tablets is assessed using a combination of simple failure and erosion tests, and more sophisticated engineering tests. The simpler tests are often used for quality control purposes, whereas the more complex tests are used during the design of the formulation and manufacturing process in the research and development phase. Standards for tablet properties are published in the various international pharmacopeias (USP/NF, EP, JP, etc.). The hardness of tablets is the principal measure of mechanical strength. Hardness is tested using a tablet hardness tester. The units for hardness have evolved since the 1930s but are commonly measured in kilograms per square centimetre. Models of testers include the Monsanto (or Stokes) Hardness Tester from 1930, the Pfizer Hardness Tester from 1950, the Strong Cob Hardness Tester and the Heberlain (or Schleeniger) Hardness Tester.
Lubricants prevent ingredients from clumping together and from sticking to the tablet punches or capsule filling machine. Lubricants also ensure that tablet formation and ejection can occur with low friction between the solid and die wall, as well as between granules, which helps in uniform filling of the die.
Common minerals like talc or silica, and fats, e.g. vegetable stearin, magnesium stearate or stearic acid are the most frequently used lubricants in tablets or hard gelatin capsules.
Manufacturing
Manufacture of the tableting blend
In the tablet pressing process, the appropriate amount of active ingredient must be in each tablet. Hence, all the ingredients should be well mixed. If a sufficiently homogenous mix of the components cannot be obtained with simple blending processes, the ingredients must be granulated prior to compression to assure an even distribution of the active compound in the final tablet. Two basic techniques are used to granulate powders for compression into tablets: wet granulation and dry granulation. Powders that can be mixed well do not require granulation and can be compressed into tablets through direct compression ("DC"). Direct compression is desirable as it is quicker. There is less processing, equipment, labor, and energy consumption. However, DC is difficult when a formulation has a high content of poorly compressible active ingredients.
Wet granulation
Wet granulation is a process of using a liquid binder to lightly agglomerate the powder mixture. The amount of liquid has to be properly controlled, as over-wetting will cause the granules to be too hard and under-wetting will cause them to be too soft and friable. Aqueous solutions have the advantage of being safer to deal with than solvent-based systems but may not be suitable for drugs which are degraded by hydrolysis.
Procedure
The active ingredient and excipients are weighed and mixed.
The wet granulate is prepared by adding the liquid binder–adhesive to the powder blend and mixing thoroughly. Examples of binders/adhesives include aqueous preparations of cornstarch, natural gums such as acacia, cellulose derivatives such as methyl cellulose, gelatin, and povidone.
Screening the damp mass through a mesh to form pellets or granules.
Drying the granulation. A conventional tray-dryer or fluid-bed dryer are most commonly used.
After the granules are dried, they are passed through a screen of smaller size than the one used for the wet mass to create granules of uniform size.
Low shear wet granulation processes use very simple mixing equipment, and can take a considerable time to achieve a uniformly mixed state. High shear wet granulation processes use equipment that mixes the powder and liquid at a very fast rate, and thus speeds up the manufacturing process. Fluid bed granulation is a multiple-step wet granulation process performed in the same vessel to pre-heat, granulate, and dry the powders. It is used because it allows close control of the granulation process.
Dry granulation
Dry granulation processes create granules by light compaction of the powder blend under low pressures. The compacts so-formed are broken up gently to produce granules (agglomerates). This process is often used when the product to be granulated is sensitive to moisture and heat. Dry granulation can be conducted on a tablet press using slugging tooling or on a roll press called a roller compactor. Dry granulation equipment offers a wide range of pressures to attain proper densification and granule formation. Dry granulation is simpler than wet granulation, therefore the cost is reduced. However, dry granulation often produces a higher percentage of fine granules, which can compromise the quality or create yield problems for the tablet. Dry granulation requires drugs or excipients with cohesive properties, and a 'dry binder' may need to be added to the formulation to facilitate the formation of granules.
Hot melt extrusion
Hot melt extrusion is utilized in pharmaceutical solid oral dose processing to enable delivery of drugs with poor solubility and bioavailability. Hot melt extrusion has been shown to molecularly disperse poorly soluble drugs in a polymer carrier increasing dissolution rates and bioavailability. The process involves the application of heat, pressure and agitation to mix materials together and 'extrude' them through a die. Twin-screw high shear extruders blend materials and simultaneously break up particles. The extruded particles can then be blended and compressed into tablets or filled into capsules.
Granule lubrication
After granulation, a final lubrication step is used to ensure that the tableting blend does not stick to the equipment during the tableting process. This usually involves low shear blending of the granules with a powdered lubricant, such as magnesium stearate or stearic acid.
Manufacture of the tablets
Whatever process is used to make the tableting blend, the process of making a tablet by powder compaction is very similar. First, the powder is filled into the die from above. The mass of powder is determined by the position of the lower punch in the die, the cross-sectional area of the die, and the powder density. At this stage, adjustments to the tablet weight are normally made by repositioning the lower punch. After die filling, the upper punch is lowered into the die and the powder is uniaxially compressed to a porosity of between 5 and 20%. The compression can take place in one or two stages (main compression, and, sometimes, pre-compression or tamping) and for commercial production occurs very fast (500–50 mg per tablet). Finally, the upper punch is pulled up and out of the die (decompression), and the tablet is ejected from the die by lifting the lower punch until its upper surface is flush with the top face of the die. This process is repeated for each tablet.
Common problems encountered during tablet manufacturing operations include:
Fluctuations in tablet weight, usually caused by uneven powder flow into the die due to poor powder flow properties.
Fluctuations in dosage of the Active Pharmaceutical Ingredient, caused by uneven distribution of the API in the tableting blend (either due to poor mixing or separation in process).
Sticking of the powder blend to the tablet tooling, due to inadequate lubrication, worn or dirty tooling, or a sticky powder formulation
Capping, lamination or chipping. This is caused by air being compressed with the tablet formulation and then expanding when the punch is released: if this breaks the tablet apart, it can be due to incorrect machine settings, or due to incorrect formulation: either because the tablet formulation is too brittle or not adhesive enough, or because the powder being fed to the tablet press contains too much air (has too low bulk density).
Capping can also occur due to high moisture content.
Consequently, permanent consistency checks are required during the manufacturing process.
Tablet compaction simulator
Tablet formulations are designed and tested using a laboratory machine called a Tablet Compaction Simulator or Powder Compaction Simulator. This is a computer controlled device that can measure the punch positions, punch pressures, friction forces, die wall pressures, and sometimes the tablet internal temperature during the compaction event. Numerous experiments with small quantities of different mixtures can be performed to optimise a formulation. Mathematically corrected punch motions can be programmed to simulate any type and model of production tablet press. Initial quantities of active pharmaceutical ingredients are very expensive to produce, and using a Compaction Simulator reduces the amount of powder required for product development.
Tablet presses
Tablet presses, also called tableting machines, range from small, inexpensive bench-top models that make one tablet at a time (single-station presses), with only around a half-ton pressure, to large, computerized, industrial models (multi-station rotary presses) that can make hundreds of thousands to millions of tablets an hour with much greater pressure. The tablet press is an essential piece of machinery for any pharmaceutical and nutraceutical manufacturer. Tablet presses must allow the operator to adjust the position of the lower and upper punches accurately, so that the tablet weight, thickness and density/hardness can each be controlled. This is achieved using a series of cams, rollers, or tracks that act on the tablet tooling (punches). Mechanical systems are also incorporated for die filling, and for ejecting and removing the tablets from the press after compression. Pharmaceutical tablet presses are required to be easy to clean and quick to reconfigure with different tooling, because they are usually used to manufacture many different products.
There are two main standards of tablet tooling used in pharmaceutical industry: American standard TSM and European standard EU. TSM and EU configurations are similar to each other but cannot be interchanged.
Modern tablet presses reach output volumes of up to 1,700,000 tablets per hour. These huge volumes require frequent in-process quality control for the tablet weight, thickness and hardness. Due to efforts to reduce rejects rates and machine down-time, automated tablet testing devices are used on-line with the tablet press or off-line in the IPC-labs.
Tablet coating
Many tablets today are coated after being pressed. Although sugar-coating was popular in the past, the process has many drawbacks. Modern tablet coatings are polymer and polysaccharide based, with plasticizers and pigments included. Tablet coatings must be stable and strong enough to survive the handling of the tablet, must not make tablets stick together during the coating process, and must follow the fine contours of embossed characters or logos on tablets. Coatings are necessary for tablets that have an unpleasant taste, and a smoother finish makes large tablets easier to swallow. Tablet coatings are also useful to extend the shelf-life of components that are sensitive to moisture or oxidation. Special coatings (for example with pearlescent effects) can enhance brand recognition.
If the active ingredient of a tablet is sensitive to acid, or is irritant to the stomach lining, an enteric coating can be used, which is resistant to stomach acid, and dissolves in the less acidic area of the intestines. Enteric coatings are also used for medicines that can be negatively affected by taking a long time to reach the small intestine, where they are absorbed. Coatings are often chosen to control the rate of dissolution of the drug in the gastrointestinal tract. Some drugs are absorbed better in certain parts of the digestive system. If this part is the stomach, a coating is selected that dissolves quickly and easily in acid. If the rate of absorption is best in the large intestine or colon, a coating is used that is acid resistant and dissolves slowly to ensure that the tablet reaches that point before dispersing. To measure the disintegration time of the tablet coating and the tablet core, automatic disintegration testers are used which are able to determine the complete disintegration process of a tablet by measuring the rest height of the thickness with every upward stroke of the disintegration tester basket.
There are two types of coating machines used in the pharmaceutical industry: coating pans and automatic coaters. Coating pans are used mostly to sugar coat pellets. Automatic coaters are used for all kinds of coatings; they can be equipped with a remote control panel, a dehumidifier, and dust collectors. An explosion-proof design is required for applying coatings that contain alcohol.
Pill-splitters
It is sometimes necessary to split tablets into halves or quarters. Tablets are easier to break accurately if scored, but there are devices called pill-splitters which cut unscored and scored tablets. Tablets with special coatings (for example, enteric coatings or controlled-release coatings) should not be broken before use, as this exposes the tablet core to the digestive juices, circumventing the intended delayed-release effect.
| Biology and health sciences | General concepts_2 | Health |
78011 | https://en.wikipedia.org/wiki/EDVAC | EDVAC | EDVAC (Electronic Discrete Variable Automatic Computer) was one of the earliest electronic computers. It was built by Moore School of Electrical Engineering at the University of Pennsylvania. Along with ORDVAC, it was a successor to the ENIAC. Unlike ENIAC, it was binary rather than decimal, and was designed to be a stored-program computer.
ENIAC inventors, John Mauchly and J. Presper Eckert, proposed the EDVAC's construction in August 1945. A contract to build the new computer was signed in April 1946 with an initial budget of US$100,000. EDVAC was delivered to the Ballistic Research Laboratory in 1949. The Ballistic Research Laboratory became a part of the US Army Research Laboratory in 1952.
Functionally, EDVAC was a binary serial computer with automatic addition, subtraction, multiplication, programmed division and automatic checking with an ultrasonic serial memory having a capacity of 1,024 44-bit words. EDVAC's average addition time was 864 microseconds and its average multiplication time was 2,900 microseconds.
Project and plan
ENIAC inventors John Mauchly and J. Presper Eckert proposed EDVAC's construction in August 1944, and design work for EDVAC commenced before ENIAC was fully operational. The design would implement a number of important architectural and logical improvements conceived during the ENIAC's construction and would incorporate a high-speed serial-access memory. Like the ENIAC, the EDVAC was built for the U.S. Army's Ballistics Research Laboratory at the Aberdeen Proving Ground by the University of Pennsylvania's Moore School of Electrical Engineering. Eckert and Mauchly and the other ENIAC designers were joined by John von Neumann in a consulting role; von Neumann summarized and discussed logical design developments in the 1945 First Draft of a Report on the EDVAC, written between February and June of that year. Later in September 1945, Eckert and Mauchly followed up with a progress report on automatic high-speed computing for the EDVAC.
In early 1946, months after the completion of ENIAC, the University of Pennsylvania adopted a new patent policy, which would have required Eckert and Mauchly to assign all their patents to the university if they stayed beyond spring of that year. Unable to reach an agreement with the university, the duo left the Moore School of Electrical Engineering in March 1946, along with many of the senior engineering staff. Simultaneously, the duo founded the Electronic Control Company (later renamed the Eckert-Mauchly Computer Corporation) in Philadelphia. A contract to build the new computer was signed in April 1946 with an initial budget of US$100,000. Later in August of that year, during the last of the Moore School Lectures, the Moore School team members were proposing new technological designs for the computer and its stored program concept. The contract named the device the Electronic Discrete Variable Automatic Calculator. The final cost of EDVAC, however, was similar to the ENIAC's, at just under $500,000.
The Raytheon Company was a subcontractor on EDVAC machines.
Technical description
The EDVAC was a binary serial computer with automatic addition, subtraction, multiplication, programmed division and automatic checking with an ultrasonic serial memory capacity of 1,024 44-bit words, thus giving a memory, in modern terms, of 5.6 kilobytes.
Physically, the computer comprised the following components:
a magnetic tape reader-recorder (Wilkes 1956:36 describes this as a wire recorder.)
a control unit with an oscilloscope
a dispatcher unit to receive instructions from the control and memory and direct them to other units
a computational unit to perform arithmetic operations on a pair of numbers and send the result to memory after checking on a duplicate unit
a timer
a dual memory unit consisting of two sets of 64 mercury acoustic delay lines of eight words capacity on each line
three temporary delay-line tanks each holding a single word
EDVAC's average addition time was 864 microseconds (about 1,160 operations per second) and its average multiplication time was 2,900 microseconds (about 340 operations per second). Time for an operation depended on memory access time, which varied depending on the memory address and the current point in the serial memory's recirculation cycle.
The computer had 5,937 vacuum tubes and 12,000 diodes, and consumed 56 kW of power. It covered 490 ft² (45.5 m2) of floor space and weighed . The full complement of operating personnel was thirty people per eight-hour shift.
EDVAC could also do floating-point arithmetic. It used 33 bits for the mantissa and one bit for its sign. It used 10 bits for the power of 2, including the sign bit.
For executable instructions, the 44-bit word was divided into four 10-bit addresses and four bits to encode the index of an operation. The first two addresses were to the numbers in memory being used in the operation, the third address was for the memory location to store the result, and the fourth address was the location of the next instruction to be executed. Only 12 of the possible 16 instructions were used.
Impact on future computer design
John Von Neumann's famous EDVAC monograph, First Draft of a Report on the EDVAC, proposed the main enhancement to its design that embodied the principal "stored-program" concept that we now call the Von Neumann architecture. This was the storing of the program in the same memory as the data. The British computers EDSAC at Cambridge and the Manchester Baby were the first working computers that followed this design, and it has been followed by the great majority of computers made since. Having the program and data in different memories is now called the Harvard architecture to distinguish it.
Installation and operation
EDVAC was delivered to the Ballistics Research Laboratory in 1949. After a number of problems had been discovered and solved, the computer began operation in 1951 although only on a limited basis.
In 1952 (April/May), it was running over hours a day (period from 15 April to 31 May, used for 341 hours).
By 1957, EDVAC was running over 20 hours a day with error-free run time averaging 8 hours. EDVAC received a number of upgrades including punch-card I/O in 1954, extra memory in slower magnetic drum form in 1955, and a floating-point arithmetic unit in 1958.
EDVAC ran until 1962 when it was replaced by BRLESC.
| Technology | Early computers | null |
78027 | https://en.wikipedia.org/wiki/Whirlwind%20I | Whirlwind I | Whirlwind I was a Cold War-era vacuum-tube computer developed by the MIT Servomechanisms Laboratory for the U.S. Navy. Operational in 1951, it was among the first digital electronic computers that operated in real-time for output, and the first that was not simply an electronic replacement of older mechanical systems.
It was one of the first computers to calculate in bit-parallel (rather than bit-serial), and was the first to use magnetic-core memory.
Its development led directly to the Whirlwind II design used as the basis for the United States Air Force SAGE air defense system, and indirectly to almost all business computers and minicomputers in the 1960s, particularly because of the mantra "short word length, speed, people."
Background
During World War II, the U.S. Navy's Naval Research Lab approached MIT about the possibility of creating a computer to drive a flight simulator for training bomber crews. They envisioned a fairly simple system in which the computer would continually update a simulated instrument panel based on control inputs from the pilots. Unlike older systems such as the Link Trainer, the system they envisioned would have a considerably more realistic aerodynamics model that could be adapted to any type of plane. This was an important consideration at the time, when many new designs were being introduced into service.
The Servomechanisms Lab in MIT building 32 conducted a short survey that concluded such a system was possible. The Navy's Office of Naval Research decided to fund development under Project Whirlwind (and its sister projects, Project Typhoon and Project Cyclone, with other institutions), and the lab placed Jay Forrester in charge of the project. They soon built a large analog computer for the task, but found that it was inaccurate and inflexible. Solving these problems in a general way would require a much larger system, perhaps one so large as to be impossible to construct. Judy Clapp was an early senior technical member of this team.
Perry Crawford, another member of the MIT team, saw a demonstration of ENIAC in 1945. He then suggested that a digital computer would be the best solution. Such a machine would allow the accuracy of simulations to be improved with the addition of more code in the computer program, as opposed to adding parts to the machine. As long as the machine was fast enough, there was no theoretical limit to the complexity of the simulation.
Until this point, all computers constructed were dedicated to single tasks, and run in batch mode. A series of inputs were set up in advance and fed into the computer, which would work out the answers and print them. This was not appropriate for the Whirlwind system, which needed to operate continually on an ever-changing series of inputs. Speed became a major issue: whereas with other systems it simply meant waiting longer for the printout, with Whirlwind it meant seriously limiting the amount of complexity the simulation could include.
Technical description
Design and construction
By 1947, Forrester and collaborator Robert Everett completed the design of a high-speed stored-program computer for this task. Most computers of the era operated in bit-serial mode, using single-bit arithmetic and feeding in large words, often 48 or 60 bits in size, one bit at a time. This was simply not fast enough for their purposes, so Whirlwind included sixteen such math units, operating on a complete 16-bit word every cycle in bit-parallel mode. Ignoring memory speed, Whirlwind ("20,000 single-address operations per second" in 1951) was essentially sixteen times as fast as other machines. Today, almost all CPUs perform arithmetic in "bit-parallel" mode.
The word size was selected after some deliberation. The machine worked by passing in a single address with almost every instruction, thereby reducing the number of memory accesses. For operations with two operands, adding for instance, the "other" operand was assumed to be the last one loaded. Whirlwind operated much like a reverse Polish notation calculator in this respect; except there was no operand stack, only an accumulator. The designers felt that 2048 words of memory would be the minimum usable amount, requiring 11 bits to represent an address, and that 16 to 32 instructions would be the minimum for another five bits — and so it was 16 bits.
The Whirlwind design incorporated a control store driven by a master clock. Each step of the clock selected one or more signal lines in a diode matrix that enabled gates and other circuits on the machine. A special switch directed signals to different parts of the matrix to implement different instructions. In the early 1950s, Whirlwind I "would crash every 20 minutes on average."
Whirlwind construction started in 1948, an effort that employed 175 people, including 70 engineers and technicians. The use of carry save multiplication appears to have been first introduced in the Whirlwind computer in the late 1940s. In the third quarter of 1949, the computer was advanced enough to solve an equation and display its solution on an oscilloscope, and even for the first animated and interactive computer graphic game. Finally Whirlwind "successfully accomplished digital computation of interception courses" on April 20, 1951. The project's budget was approximately $1 million a year, which was vastly higher than the development costs of most other computers of the era. After three years, the Navy had lost interest. However, during this time the Air Force had become interested in using computers to help the task of ground controlled interception, and the Whirlwind was the only machine suitable to the task. They took up development under Project Claude.
Whirlwind weighed and occupied over .
The memory subsystem
The original machine design called for 2048 (2K) words of 16 bits each of random-access storage. The only two available memory technologies in 1949 that could hold this much data were mercury delay lines and electrostatic storage.
A mercury delay line consisted of a long tube filled with mercury, a mechanical transducer on one end, and a microphone on the other end, much like a spring reverb unit later used in audio processing. Pulses were sent into the mercury delay line at one end, and took a certain amount of time to reach the other end. They were detected by the microphone, amplified, reshaped into the correct pulse shape, and sent back into the delay line. Thus, the memory was said to recirculate.
Mercury delay lines operated at about the speed of sound, so were very slow in computer terms, even by the standards of the computers of the late 1940s and 1950s. The speed of sound in mercury was also very dependent on temperature. Since a delay line held a defined number of bits, the frequency of the clock had to change with the temperature of the mercury. If there were many delay lines and they did not all have the same temperature at all times, the memory data could easily become corrupted.
The Whirlwind designers quickly discarded the delay line as a possible memory—it was both too slow for the envisioned flight simulator, and too unreliable for a reproducible production system, for which Whirlwind was intended to be a functional prototype.
The alternative form of memory was known as "electrostatic". This was a cathode ray tube memory, similar in many aspects to an early TV picture tube or oscilloscope tube. An electron gun sent a beam of electrons to the far end of the tube, where they impacted a screen. The beam would be deflected to land at a particular spot on the screen. The beam could then build up a negative charge at that point, or change a charge that was already there. By measuring the beam current it could be determined whether the spot was originally a zero or a one, and a new value could be stored by the beam.
There were several forms of electrostatic memory tubes in existence in 1949. The best known today is the Williams tube, developed in England, but there were a number of others that had been developed independently by various research labs. The Whirlwind engineers considered the Williams tube, but determined that the dynamic nature of the storage and the need for frequent refresh cycles was incompatible with the design goals for Whirlwind I. Instead, they settled on a design that was being developed at the MIT Radiation Laboratory. This was a dual-gun electron tube. One gun produced a sharply-focused beam to read or write individual bits. The other gun was a "flood gun" that sprayed the entire screen with low-energy electrons. As a result of the design, this tube was more of a static RAM that did not require refresh cycles, unlike the dynamic RAM Williams tube.
In the end the choice of this tube was unfortunate. The Williams tube was considerably better developed, and despite the need for refresh could easily hold 1024 bits per tube, and was quite reliable when operated correctly. The MIT tube was still in development, and while the goal was to hold 1024 bits per tube, this goal was never reached, even several years after the plan had called for full-size functional tubes. Also, the specifications had called for an access time of six microseconds, but the actual access time was around 30 microseconds. Since the basic cycle time of the Whirlwind I processor was determined by the memory access time, the entire processor was slower than designed.
Magnetic-core memory
Jay Forrester was desperate to find a suitable memory replacement for his computer. Initially the computer only had 32 words of storage, and 27 of these words were read-only registers made of toggle switches. The remaining five registers were flip-flop storage, with each of the five registers being made from more than 30 vacuum tubes. This "test storage", as it was known, was intended to allow checkout of the processing elements while the main memory was not ready. The main memory was so late that the first experiments of tracking airplanes with live radar data were done using a program manually set into test storage. Forrester came across an advertisement for a new magnetic material being produced by a company. Recognizing that this had the potential to be a data storage medium, Forrester obtained a workbench in the corner of the lab, and got several samples of the material to experiment with. Then for several months he spent as much time in the lab as he did in the office managing the entire project.
At the end of those months, he had invented the basics of magnetic-core memory and demonstrated that it was likely to be feasible. His demonstration consisted of a small core plane of 32 cores, each three-eighths of an inch in diameter. Having demonstrated that the concept was practical, it needed only to be reduced to a workable design. In the fall of 1949, Forrester enlisted graduate student William N. Papian to test dozens of individual cores, to determine those with the best properties. Papian's work was bolstered when Forrester asked student Dudley Allen Buck to work on the material and assigned him to the workbench, while Forrester went back to full-time project management. (Buck would go on to invent the cryotron and content-addressable memory at the lab.)
After approximately two years of further research and development, they were able to demonstrate a core plane that was made of 32 by 32, or 1024 cores, holding 1024 bits of data. Thus, they had reached the originally intended storage size of an electrostatic tube, a goal that had not yet been reached by the tubes themselves, only holding 512 bits per tube in the latest design generation. Very quickly, a 1024-word core memory was fabricated, replacing the electrostatic memory. The electrostatic memory design and production was summarily canceled, saving a good deal of money to be reallocated to other research areas. Two additional core memory units were later fabricated, increasing the total memory size available.
Vacuum tubes
The design used approximately 5,000 vacuum tubes.
The large number of tubes used in Whirlwind resulted in a problematic failure rate since a single tube failure could cause a system failure. The standard pentode at the time was the 6AG7, but testing in 1948 determined that its expected lifetime in service was too short for this application. Consequently, the 7AD7 was chosen instead, but this also had too high a failure rate in service. An investigation into the cause of the failures found that silicon in the tungsten alloy of the heater filament was causing cathode poisoning; deposits of barium orthosilicate forming on the cathode reduce or prevent its function of emitting electrons. The 7AK7 tube with a high-purity tungsten filament was then specially developed for Whirlwind by Sylvania.
Cathode poisoning is at its worst when the tube is being run in cut-off with the heater on. Commercial tubes were intended for radio (and later, television) applications where they are rarely run in this state. Analog applications like these keep the tube in the linear region, whereas digital applications switch the tube between cut-off and full conduction, passing only briefly through the linear region. Further, commercial manufacturers expected their tubes to only be in use for a few hours per day. To ameliorate this issue, the heaters were turned off on valves not expected to switch for long periods. The heater voltage was turned on and off with a slow ramp waveform to avoid thermal shock to the heater filaments.
Even these measures were not enough to achieve the required reliability. Incipient faults were proactively sought by testing the valves during maintenance periods. They were subject to stress tests called marginal testing because they applied voltages and signals to the valves right up to their design margins. These tests were designed to bring on early failure of valves that would otherwise have failed while in service. They were carried out automatically by a test program. The maintenance statistics for 1950 show the success of these measures. Of the 1,622 7AD7 tubes in use, 243 failed, of which 168 were found by marginal testing. Of the 1,412 7AK7 tubes in use, 18 failed, of which only 2 failed during marginal checking. As a result, Whirlwind was far more reliable than any commercially available machine.
Many other features of the Whirlwind tube testing regime were not standard tests and required specially built equipment. One condition that required special testing was momentary shorting on a few tubes caused by small objects like lint inside the tube. Occasional spurious short pulses are a minor problem, or even entirely unnoticeable, in analog circuits, but are likely to be disastrous in a digital circuit. These did not show up on standard tests but could be discovered manually by tapping the glass envelope. A thyratron-triggered circuit was built to automate this test.
Air defense networks
After connection to the experimental Microwave Early Warning (MEW) radar at Hanscom Field using Jack Harrington's equipment and commercial phone lines, aircraft were tracked by Whirlwind I. The Cape Cod System subsequently demonstrated computerized air defence covering southern New England. Signals from three long range (AN/FPS-3) radars, eleven gap-filler radars, and three height-finding radars were transmitted over telephone lines to the Whirlwind I computer in Cambridge, Massachusetts. The Whirlwind II design for a larger and faster machine (never completed) was the basis for the SAGE air defense system IBM AN/FSQ-7 Combat Direction Central.
Legacy
The Whirlwind used approximately 5,000 vacuum tubes. An effort was also started to convert the Whirlwind design to a transistorized form, led by Ken Olsen and known as the TX-0. TX-0 was very successful and plans were made to make an even larger version known as TX-1. However this project was far too ambitious and had to be scaled back to a smaller version known as TX-2. Even this version proved troublesome, and Olsen left in mid-project to start Digital Equipment Corporation (DEC). DEC's PDP-1 was essentially a collection of TX-0 and TX-2 concepts in a smaller package.
After supporting SAGE, Whirlwind I was rented ($1/yr) from June 30, 1959, until 1974 by project member, William M. Wolf (1928-2015). The power to run the machine cost $2500 per month, and the Wolf Research and Development Corporation did work for the Air Force and Buckminster Fuller's World Game. Ultimately moving Whirlwind I cost $250,000 and the company made $100,000 on it. Wolf R&D Corporation was sold to EG&G in 1967 for $5.5 million.
Ken Olsen and Robert Everett saved the machine, which became the basis for the Boston Computer Museum in 1979. Although much of the machine was lost when decommissioned, many of its components are now in the collection of the Computer History Museum in Mountain View, California and the MIT Museum.
As of February 2009, a core memory unit is displayed at the Charles River Museum of Industry & Innovation in Waltham, Massachusetts. One plane of core memory, on loan from the Computer History Museum, is on shown as part of the Historic Computer Science displays at the Gates Computer Science Building, Stanford.
The building which housed Whirlwind was until recently home to MIT's campus-wide IT department, Information Services & Technology and in 1997–1998, it was restored to its original exterior design.
| Technology | Early computers | null |
78130 | https://en.wikipedia.org/wiki/Max-flow%20min-cut%20theorem | Max-flow min-cut theorem | In computer science and optimization theory, the max-flow min-cut theorem states that in a flow network, the maximum amount of flow passing from the source to the sink is equal to the total weight of the edges in a minimum cut, i.e., the smallest total weight of the edges which if removed would disconnect the source from the sink.
This is a special case of the duality theorem for linear programs and can be used to derive Menger's theorem and the Kőnig–Egerváry theorem.
Definitions and statement
The theorem equates two quantities: the maximum flow through a network, and the minimum capacity of a cut of the network. To state the theorem, each of these notions must first be defined.
Network
A network consists of
a finite directed graph , where V denotes the finite set of vertices and is the set of directed edges;
a source and a sink ;
a capacity function, which is a mapping denoted by or for . It represents the maximum amount of flow that can pass through an edge.
Flows
A flow through a network is a mapping denoted by or , subject to the following two constraints:
Capacity Constraint: For every edge ,
Conservation of Flows: For each vertex apart from and (i.e. the source and sink, respectively), the following equality holds:
A flow can be visualized as a physical flow of a fluid through the network, following the direction of each edge. The capacity constraint then says that the volume flowing through each edge per unit time is less than or equal to the maximum capacity of the edge, and the conservation constraint says that the amount that flows into each vertex equals the amount flowing out of each vertex, apart from the source and sink vertices.
The value of a flow is defined by
where as above is the source and is the sink of the network. In the fluid analogy, it represents the amount of fluid entering the network at the source. Because of the conservation axiom for flows, this is the same as the amount of flow leaving the network at the sink.
The maximum flow problem asks for the largest flow on a given network.
Maximum Flow Problem. Maximize , that is, to route as much flow as possible from to .
Cuts
The other half of the max-flow min-cut theorem refers to a different aspect of a network: the collection of cuts. An s-t cut is a partition of such that and . That is, an s-t cut is a division of the vertices of the network into two parts, with the source in one part and the sink in the other. The cut-set of a cut is the set of edges that connect the source part of the cut to the sink part:
Thus, if all the edges in the cut-set of are removed, then no positive flow is possible, because there is no path in the resulting graph from the source to the sink.
The capacity of an s-t cut is the sum of the capacities of the edges in its cut-set,
where if and , otherwise.
There are typically many cuts in a graph, but cuts with smaller weights are often more difficult to find.
Minimum s-t Cut Problem. Minimize , that is, determine and such that the capacity of the s-t cut is minimal.
Main theorem
In the above situation, one can prove that the value of any flow through a network is less than or equal to the capacity of any s-t cut, and that furthermore a flow with maximal value and a cut with minimal capacity exist. The main theorem links the maximum flow value with the minimum cut capacity of the network.
Max-flow min-cut theorem. The maximum value of an s-t flow is equal to the minimum capacity over all s-t cuts.
Example
The figure on the right shows a flow in a network. The numerical annotation on each arrow, in the form f/c, indicates the flow (f) and the capacity (c) of the arrow. The flows emanating from the source total five (2+3=5), as do the flows into the sink (2+3=5), establishing that the flow's value is 5.
One s-t cut with value 5 is given by S={s,p} and T={o, q, r, t}. The capacities of the edges that cross this cut are 3 and 2, giving a cut capacity of 3+2=5. (The arrow from o to p is not considered, as it points from T back to S.)
The value of the flow is equal to the capacity of the cut, showing that the flow is a maximal flow and the cut is a minimal cut.
Note that the flow through each of the two arrows that connect S to T is at full capacity; this is always the case: a minimal cut represents a 'bottleneck' of the system.
Linear program formulation
The max-flow problem and min-cut problem can be formulated as two primal-dual linear programs.
The max-flow LP is straightforward. The dual LP is obtained using the algorithm described in dual linear program: the variables and sign constraints of the dual correspond to the constraints of the primal, and the constraints of the dual correspond to the variables and sign constraints of the primal. The resulting LP requires some explanation. The interpretation of the variables in the min-cut LP is:
The minimization objective sums the capacity over all the edges that are contained in the cut.
The constraints guarantee that the variables indeed represent a legal cut:
The constraints (equivalent to ) guarantee that, for non-terminal nodes u,v, if u is in S and v is in T, then the edge (u,v) is counted in the cut ().
The constraints (equivalent to ) guarantee that, if v is in T, then the edge (s,v) is counted in the cut (since s is by definition in S).
The constraints (equivalent to ) guarantee that, if u is in S, then the edge (u,t) is counted in the cut (since t is by definition in T).
Note that, since this is a minimization problem, we do not have to guarantee that an edge is not in the cut - we only have to guarantee that each edge that should be in the cut, is summed in the objective function.
The equality in the max-flow min-cut theorem follows from the strong duality theorem in linear programming, which states that if the primal program has an optimal solution, x*, then the dual program also has an optimal solution, y*, such that the optimal values formed by the two solutions are equal.
Application
Cederbaum's maximum flow theorem
The maximum flow problem can be formulated as the maximization of the electrical current through a network composed of nonlinear resistive elements. In this formulation, the limit of the current between the input terminals of the electrical network as the input voltage approaches , is equal to the weight of the minimum-weight cut set.
Generalized max-flow min-cut theorem
In addition to edge capacity, consider there is capacity at each vertex, that is, a mapping denoted by , such that the flow has to satisfy not only the capacity constraint and the conservation of flows, but also the vertex capacity constraint
In other words, the amount of flow passing through a vertex cannot exceed its capacity. Define an s-t cut to be the set of vertices and edges such that for any path from s to t, the path contains a member of the cut. In this case, the capacity of the cut is the sum of the capacity of each edge and vertex in it.
In this new definition, the generalized max-flow min-cut theorem states that the maximum value of an s-t flow is equal to the minimum capacity of an s-t cut in the new sense.
Menger's theorem
In the undirected edge-disjoint paths problem, we are given an undirected graph and two vertices and , and we have to find the maximum number of edge-disjoint s-t paths in .
Menger's theorem states that the maximum number of edge-disjoint s-t paths in an undirected graph is equal to the minimum number of edges in an s-t cut-set.
Project selection problem
In the project selection problem, there are projects and machines. Each project yields revenue and each machine costs to purchase. We want to select a subset of the project, and purchase a subset of the machines, to maximize the total profit (revenue of the selected projects minus cost of the purchased machines). We must obey the following constraint: each project specifies a set of machines which must be purchased if the project is selected. (Each machine, once purchased, can be used by any selected project.)
To solve the problem, let be the set of projects not selected and be the set of machines purchased, then the problem can be formulated as,
Since the first term does not depend on the choice of and , this maximization problem can be formulated as a minimization problem instead, that is,
The above minimization problem can then be formulated as a minimum-cut problem by constructing a network, where the source is connected to the projects with capacity , and the sink is connected by the machines with capacity . An edge with infinite capacity is added if project requires machine . The s-t cut-set represents the projects and machines in and respectively. By the max-flow min-cut theorem, one can solve the problem as a maximum flow problem.
The figure on the right gives a network formulation of the following project selection problem:
The minimum capacity of an s-t cut is 250 and the sum of the revenue of each project is 450; therefore the maximum profit g is 450 − 250 = 200, by selecting projects and .
The idea here is to 'flow' each project's profits through the 'pipes' of its machines. If we cannot fill the pipe from a machine, the machine's return is less than its cost, and the min cut algorithm will find it cheaper to cut the project's profit edge instead of the machine's cost edge.
Image segmentation problem
In the image segmentation problem, there are pixels. Each pixel can be assigned a foreground value or a background value . There is a penalty of if pixels are adjacent and have different assignments. The problem is to assign pixels to foreground or background such that the sum of their values minus the penalties is maximum.
Let be the set of pixels assigned to foreground and be the set of points assigned to background, then the problem can be formulated as,
This maximization problem can be formulated as a minimization problem instead, that is,
The above minimization problem can be formulated as a minimum-cut problem by constructing a network where the source (orange node) is connected to all the pixels with capacity , and the sink (purple node) is connected by all the pixels with capacity . Two edges () and () with capacity are added between two adjacent pixels. The s-t cut-set then represents the pixels assigned to the foreground in and pixels assigned to background in .
History
An account of the discovery of the theorem was given by Ford and Fulkerson in 1962:
"Determining a maximal steady state flow from one point to another in a network subject to capacity limitations on arcs ... was posed to the authors in the spring of 1955 by T.E. Harris, who, in conjunction with General F. S. Ross (Ret.) had formulated a simplified model of railway traffic flow, and pinpointed this particular problem as the central one suggested by the model. It was not long after this until the main result, Theorem 5.1, which we call the max-flow min-cut theorem, was conjectured and established. A number of proofs have since appeared."
Proof
Let be a network (directed graph) with and being the source and the sink of respectively.
Consider the flow computed for by Ford–Fulkerson algorithm. In the residual graph obtained for (after the final flow assignment by Ford–Fulkerson algorithm), define two subsets of vertices as follows:
: the set of vertices reachable from in
: the set of remaining vertices i.e.
Claim. , where the capacity of an s-t cut is defined by
.
Now, we know, for any subset of vertices, . Therefore, for we need:
All outgoing edges from the cut must be fully saturated.
All incoming edges to the cut must have zero flow.
To prove the above claim we consider two cases:
In , there exists an outgoing edge such that it is not saturated, i.e., . This implies, that there exists a forward edge from to in , therefore there exists a path from to in , which is a contradiction. Hence, any outgoing edge is fully saturated.
In , there exists an incoming edge such that it carries some non-zero flow, i.e., . This implies, that there exists a backward edge from to in , therefore there exists a path from to in , which is again a contradiction. Hence, any incoming edge must have zero flow.
Both of the above statements prove that the capacity of cut obtained in the above described manner is equal to the flow obtained in the network. Also, the flow was obtained by Ford-Fulkerson algorithm, so it is the max-flow of the network as well.
Also, since any flow in the network is always less than or equal to capacity of every cut possible in a network, the above described cut is also the min-cut which obtains the max-flow.
A corollary from this proof is that the maximum flow through any set of edges in a cut of a graph is equal to the minimum capacity of all previous cuts.
| Mathematics | Graph theory | null |
78160 | https://en.wikipedia.org/wiki/Cougar | Cougar | The cougar (Puma concolor) (, KOO-gər), also known as the panther, mountain lion, catamount and puma, is a large cat native to the Americas. It inhabits North, Central and South America, making it the most widely distributed wild, terrestrial mammal in the Western Hemisphere, and one of the most widespread in the world. Its range spans the Yukon, British Columbia and Alberta provinces of Canada, the Rocky Mountains and areas in the western United States. Further south, its range extends through Mexico to the Amazon Rainforest and the southern Andes Mountains in Patagonia. It is an adaptable generalist species, occurring in most American habitat types. It prefers habitats with dense underbrush and rocky areas for stalking but also lives in open areas.
The cougar is largely solitary. Its activity pattern varies from diurnality and cathemerality to crepuscularity and nocturnality between protected and non-protected areas, and is apparently correlated with the presence of other predators, prey species, livestock and humans. It is an ambush predator that pursues a wide variety of prey. Ungulates, particularly deer, are its primary prey, but it also hunts rodents. It is territorial and lives at low population densities. Individual home ranges depend on terrain, vegetation and abundance of prey. While large, it is not always the dominant apex predator in its range, yielding prey to other predators. It is reclusive and mostly avoids people. Fatal attacks on humans are rare but increased in North America as more people entered cougar habitat and built farms.
The cougar is listed as Least Concern on the IUCN Red List. Intensive hunting following European colonization of the Americas and ongoing human development into cougar habitat has caused populations to decline in most parts of its historical range. In particular, the eastern cougar population is considered to be mostly locally extinct in eastern North America since the early 20th century, with the exception of the isolated Florida panther subpopulation.
Naming and etymology
The word cougar is borrowed from the Portuguese çuçuarana, via French; it was originally derived from the Tupi language. A current form in Brazil is suçuarana. In the 17th century, Georg Marcgrave named it cuguacu ara. Marcgrave's rendering was reproduced in 1648 by his associate Willem Piso. Cuguacu ara was then adopted by John Ray in 1693. In 1774, Georges-Louis Leclerc, Comte de Buffon converted cuguacu ara to cuguar, which was later modified to "cougar" in English.
The cougar holds the Guinness record for the animal with the greatest number of names, with over 40 in English alone. "Puma" is the common name used in Latin America and most parts of Europe. The term puma is also sometimes used in the United States. The first use of puma in English dates to 1777, introduced from Spanish from the Quechua language. In the western United States and Canada, it is also called "mountain lion", a name first used in writing in 1858. Other names include "panther" (although it does not belong to the genus Panthera) and "catamount" (meaning "cat of the mountains").
Taxonomy and evolution
Felis concolor was the scientific name proposed by Carl Linnaeus in 1771 for a cat with a long tail from Brazil. The specific epithet of the name, "concolor", is Latin for "of uniform color". It was placed in the genus Puma by William Jardine in 1834. This genus is part of the Felinae. The cougar is most closely related to the jaguarundi and the cheetah.
Subspecies
Following Linnaeus's first scientific description of the cougar, 32 cougar zoological specimens were described and proposed as subspecies until the late 1980s. Genetic analysis of cougar mitochondrial DNA indicates that many of these are too similar to be recognized as distinct at a molecular level but that only six phylogeographic groups exist. The Florida panther samples showed a low microsatellite variation, possibly due to inbreeding. Following this research, the authors of Mammal Species of the World recognized the following six subspecies in 2005:
P. c. concolor includes the synonyms bangsi, incarum, osgoodi, soasoaranna, sussuarana, soderstromii, suçuaçuara, and wavula
P. c. puma includes the synonyms araucanus, concolor, patagonica, pearsoni, and puma
P. c. couguar includes arundivaga, aztecus, browni, californica, floridana, hippolestes, improcera, kaibabensis, mayensis, missoulensis, olympus, oregonensis, schorgeri, stanleyana, vancouverensis, and youngi
P. c. costaricensis
P. c. anthonyi includes acrocodia, borbensis, capricornensis, concolor, greeni, and nigra
P. c. cabrerae includes hudsonii and puma proposed by Marcelli in 1922
In 2006, the Florida panther was still referred to as a distinct subspecies P. c. coryi in research works.
, the Cat Classification Taskforce of the Cat Specialist Group recognizes only two subspecies as valid:
P. c. concolor in South America, possibly excluding the region northwest of the Andes
P. c. couguar in North and Central America and possibly northwestern South America
Evolution
The family Felidae is believed to have originated in Asia about 11 million years ago (Mya). Taxonomic research on felids remains partial, and much of what is known about their evolutionary history is based on mitochondrial DNA analysis. Significant confidence intervals exist with suggested dates. In the latest genomic study of the Felidae, the common ancestor of today's Leopardus, Lynx, Puma, Prionailurus, and Felis lineages migrated across the Bering land bridge into the Americas 8.0 to 8.5 million years ago. The lineages subsequently diverged in that order. North American felids then invaded South America 2–4Mya as part of the Great American Interchange, following the formation of the Isthmus of Panama.
The cheetah lineage is suggested by some studies to have diverged from the Puma lineage in the Americas and migrated back to Asia and Africa, while other research suggests the cheetah diverged in the Old World itself. A high level of genetic similarity has been found among North American cougar populations, suggesting they are all fairly recent descendants of a small ancestral group. Culver et al. propose the original North American cougar population was extirpated during the Pleistocene extinctions some 10,000 years ago, when other large mammals, such as Smilodon, also disappeared. North America was then repopulated by South American cougars.
A coprolite identified as from a cougar was excavated in Argentina's Catamarca Province and dated to 17,002–16,573 years old. It contained Toxascaris leonina eggs. This finding indicates that the cougar and the parasite have existed in South America since at least the Late Pleistocene. The oldest fossil record of a cougar (Puma concolor) in South America (Argentina) is a partial skull from the late Calabrian (Ensenadan) age.
Characteristics
The head of the cougar is round, and the ears are erect. Its powerful forequarters, neck, and jaw serve to grasp and hold large prey. It has four retractile claws on its hind paws and five on its forepaws, of which one is a dewclaw. The larger front feet and claws are adaptations for clutching prey.
Cougars are slender and agile members of the Felidae. They are the fourth largest cat species worldwide; adults stand about tall at the shoulders. Adult males are around long from nose to tail tip, and females average , with overall ranges between nose to tail suggested for the species in general. Of this length, the tail typically accounts for . Males generally weigh . Females typically weigh between . Cougar size is smallest close to the equator and larger towards the poles. The largest recorded cougar, shot in 1901, weighed ; claims of and have been reported, though they were probably exaggerated. Male cougars in North America average , while the average female in the same region averages about . On average, adult male cougars in British Columbia weigh and adult females , though several male cougars in British Columbia weighed between .
Depending on the locality, cougars can be smaller or bigger than jaguars but are less muscular and not as powerfully built, so on average, their weight is less. Whereas the size of cougars tends to increase as much as distance from the equator increases, which crosses the northern portion of South America, jaguars are generally smaller north of the Amazon River in South America and larger south of it. For example, while South American jaguars are comparatively large, and may exceed , North American jaguars in Mexico's Chamela-Cuixmala Biosphere Reserve weigh approximately , about the same as female cougars.
Cougar coloring is plain (hence the Latin concolor ["one color"] in the scientific name) but can vary greatly across individuals and even siblings. The coat is typically tawny, but it ranges from silvery-grey to reddish with lighter patches on the underbody, including the jaws, chin, and throat. Infants are spotted and born with blue eyes and rings on their tails; juveniles are pale, and dark spots remain on their flanks. A leucistic individual was seen in Serra dos Órgãos National Park in Rio de Janeiro in 2013 when it was recorded by a camera trap, indicating that pure white individuals do exist within the species, though they are extremely rare.
The cougar has large paws and proportionally the largest hind legs in the Felidae, allowing for great leaping and powerful short sprints. It can leap from the ground up to high into a tree.
Distribution and habitat
The cougar has the most extensive range of any wild land animal in the Americas, spanning 110 degrees of latitude from the Yukon in Canada to the southern Andes in Chile. The species was extirpated from eastern North America, aside from Florida, but they may be recolonizing their former range and isolated populations have been documented east of their contemporary ranges in both the Midwestern US and Canada.
The cougar lives in all forest types, lowland and mountainous deserts, and in open areas with little vegetation up to an elevation of . In the Santa Ana Mountains, it prefers steep canyons, escarpments, rim rocks and dense brush. In Mexico, it was recorded in the Sierra de San Carlos. In the Yucatán Peninsula, it inhabits secondary and semi-deciduous forests in El Eden Ecological Reserve. In El Salvador, it was recorded in the lower montane forest in Montecristo National Park and in a river basin in the Morazán Department above in 2019.
In Colombia, it was recorded in a palm oil plantation close to a riparian forest in the Llanos Basin, and close to water bodies in the Magdalena River Valley.
In the human-modified landscape of central Argentina, it inhabits bushland with abundant vegetation cover and prey species.
Behavior and ecology
Cougars are an important keystone species in Western Hemisphere ecosystems, linking numerous species at many trophic levels. In a comprehensive literature review of more than 160 studies on cougar ecology, ecological interactions with 485 other species in cougar-inhabited ecosystems have been shown to involve different areas of interaction, ranging from the use of other species as food sources and prey, fear effects on potential prey, effects from carcass remains left behind, to competitive effects on other predator species in shared habitat. The most common research topic in the literature used here was the cougar's diet and its prey's regulation.
Hunting and diet
The cougar is a generalist hypercarnivore. It prefers large mammals such as mule deer, white-tailed deer, elk, moose, mountain goat and bighorn sheep. It opportunistically takes smaller prey such as rodents, lagomorphs, smaller carnivores, birds, and even domestic animals, including pets. The mean weight of cougar vertebrate prey increases with its body weight and is lower in areas closer to the equator. A survey of North America research found 68% of prey items were ungulates, especially deer. Only the Florida panther showed variation, often preferring feral hogs and armadillos. Cougars have been known to prey on introduced gemsbok populations in New Mexico. One individual cougar was recorded as hunting 29 gemsbok, which made up 58% of its recorded kills. Most gemsbok kills were neonates, but some adults were also known to have been taken. Elsewhere in the southwestern United States, they have been recorded to also prey on feral horses in the Great Basin, as well as feral donkeys in the Sonoran and Mojave Deserts.
Investigations at Yellowstone National Park showed that elk and mule deer were the cougar's primary prey; the prey base is shared with the park's wolves, with which the cougar competes for resources. A study on winter kills from November to April in Alberta showed that ungulates accounted for greater than 99% of the cougar diet. Learned, individual prey recognition was observed, as some cougars rarely killed bighorn sheep, while others relied heavily on the species.
In the Central and South American cougar range area, the ratio of deer in the diet declines. Small to mid-sized mammals, including large rodents such as the capybara, are preferred. Ungulates accounted for only 35% of prey items in one survey, about half that of North America. Competition with the larger jaguar in South America has been suggested for the decline in the size of prey items. In Central or North America, the cougar and jaguar share the same prey, depending on its abundance. Other listed prey species of the cougar include mice, porcupines, American beavers, raccoons, hares, guanacoes, peccaries, vicuñas, rheas and wild turkeys. Birds and small reptiles are sometimes preyed upon in the south, but this is rarely recorded in North America. Magellanic penguins (Spheniscus magellanicus) constitute the majority of prey items in cougar diet in Patagonia's Bosques Petrificados de Jaramillo National Park and Monte León National Park.
Although capable of sprinting, the cougar is typically an ambush predator. It stalks through brush and trees, across ledges, or other covered spots, before delivering a powerful leap onto the back of its prey and a suffocating neck bite. The cougar can break the neck of some of its smaller prey with a strong bite and momentum bearing the animal to the ground. Kills are generally estimated around one large ungulate every two weeks. The period shrinks for females raising young, and may be as short as one kill every three days when cubs are nearly mature around 15 months. The cat drags a kill to a preferred spot, covers it with brush, and returns to feed over a period of days. The cougar is generally reported to not be a scavenger, but deer carcasses left exposed for study were scavenged by cougars in California, suggesting more opportunistic behavior.
Interactions with other predators
Aside from humans, no species preys upon mature cougars in the wild, although conflicts with other predators or scavengers occur. Of the large predators in Yellowstone National Park – the grizzly and black bears, gray wolf and cougar – the massive grizzly bear appears dominant, often (though not always) able to drive a gray wolf pack, black bear or cougar off their kills. One study found that grizzlies and American black bears visited 24% of cougar kills in Yellowstone and Glacier National Parks, usurping 10% of carcasses. Bears gained up to 113%, and cougars lost up to 26% of their daily energy requirements from these encounters. In Colorado and California, black bears were found to visit 48% and 77% of kills, respectively. In general, cougars are subordinate to black bears when it comes to killing, and when bears are most active, the cats take prey more frequently and spend less time feeding on each kill. Unlike several subordinate predators from other ecosystems, cougars do not appear to exploit spatial or temporal refuges to avoid competitors.
The gray wolf and the cougar compete more directly for prey, mostly in winter. Packs of wolves can steal cougars' kills, and there are some documented cases of cougars being killed by them. One report describes a large pack of seven to 11 wolves killing a female cougar and her kittens, while in nearby Sun Valley, Idaho, a 2-year-old male cougar was found dead, apparently killed by a wolf pack. Conversely, one-to-one confrontations tend to be dominated by the cat, and there are various documented accounts where wolves have been ambushed and killed, including adult male specimens. Wolves more broadly affect cougar population dynamics and distribution by dominating territory and prey opportunities, and disrupting the feline's behavior. Preliminary research in Yellowstone, for instance, has shown displacement of the cougar by wolves. One researcher in Oregon noted: "When there is a pack around, cougars are not comfortable around their kills or raising kittens [...] A lot of times a big cougar will kill a wolf, but the pack phenomenon changes the table." Both species are capable of killing mid-sized predators, such as bobcats, Canada lynxes, wolverines and coyotes, and tend to suppress their numbers. Although cougars can kill coyotes, the latter have been documented attempting to prey on cougar cubs.
The cougar and jaguar share overlapping territory in the southern portion of its range. The jaguar tends to take the larger prey where ranges overlap, reducing both the cougar's potential size and the likelihood of direct competition between the two cats. Cougars appear better than jaguars at exploiting a broader prey niche and smaller prey.
Social spacing and interactions
The cougar is a mostly solitary animal. Only mothers and kittens live in groups, with adults meeting rarely. While generally loners, cougars will reciprocally share kills and seem to organize themselves into small communities defined by the territories of dominant males. Cats within these areas socialize more frequently with each other than with outsiders.
In the vicinity of a cattle ranch in northern Mexico, cougars exhibited nocturnal activity that overlapped foremost with the activity of calves. In a nature reserve in central Mexico, the activity of cougars was crepuscular and nocturnal, overlapping largely with the activity of the nine-banded armadillo (Dasypus novemcinctus). Cougars in the montane Abra-Tanchipa Biosphere Reserve in southeastern Mexico displayed a cathemeral activity pattern. Data from 12 years of camera trapping in the Pacific slope and Talamanca Cordillera of Costa Rica showed cougars as cathemeral. Both cougars and jaguars in the Cockscomb Basin of Belize were nocturnal but avoided each other. In a protected cloud forest in the central Andes of Colombia, cougars were active from late afternoon to shortly before sunrise and sometimes during noon and early afternoon. In protected areas of the Madidi-Tambopata Landscape in Bolivia and Peru, cougars were active throughout the day but with a tendency to nocturnal activity that overlapped with the activity of main prey species. During an 8-year-long study in a modified landscape in southeastern Brazil, male cougars were primarily nocturnal, but females were active at night and day. Cougars were diurnal in the Brazilian Pantanal, but crepuscular and nocturnal in protected areas in the Cerrado, Caatinga and ecotone biomes. Cougars in the Atlantic Forest were active throughout the day but displayed peak activity during early mornings in protected areas and crepuscular and nocturnal activity in less protected areas. In central Argentina, cougars were active day and night in protected areas but were active immediately after sunset and before sunrise outside protected areas. Cougars displayed a foremost crepuscular and nocturnal activity pattern in a ranching area in southern Argentina.
Home range sizes and overall cougar abundance depend on terrain, vegetation, and prey abundance. Research suggests a lower limit of and upper limit of of home range for males. Large male home ranges of with female ranges half that size. One female adjacent to the San Andres Mountains was found with a big range of , necessitated by poor prey abundance. Research has shown cougar abundances from 0.5 animals to as many as seven per .
Male home ranges include or overlap with females but, at least where studied, not with those of other males. The home ranges of females overlap slightly. Males create scrapes composed of leaves and duff with their hind feet, and mark them with urine and sometimes feces. When males encounter each other, they vocalize and may engage in violent conflict if neither backs down.
Cougars communicate with various vocalizations. Aggressive sounds include growls, spits, snarls, and hisses. During the mating season, estrus females produce caterwauls or yowls to attract mates, and males respond with similar vocals. Mothers and offspring keep in contact with whistles, chirps, and mews.
Reproduction and life cycle
Females reach sexual maturity at the age of 18 months to three years and are in estrus for about eight days of a 23-day cycle; the gestation period is approximately 91 days. Both adult males and females may mate with multiple partners, and a female's litter can have multiple paternities. Copulation is brief but frequent. Chronic stress can result in low reproductive rates in captivity as well as in the field.
Gestation is 82–103 days long. Only females are involved in parenting. Litter size is between one and six cubs, typically two. Caves and other alcoves that offer protection are used as litter dens. Born blind, cubs are completely dependent on their mother at first and begin to be weaned at around three months of age. As they grow, they go out on forays with their mother, first visiting kill sites and, after six months, beginning to hunt small prey on their own. Kitten survival rates are just over one per litter.
Juveniles remain with their mothers for one to two years. When a female reaches estrous again, her offspring must disperse or the male will kill them. Males tend to disperse further than females. One study has shown a high mortality rate among cougars that travel farthest from their maternal range, often due to conflicts with other cougars. In a study area in New Mexico, males dispersed farther than females, traversed large expanses of non-cougar habitat and were probably most responsible for nuclear gene flow between habitat patches.
Life expectancy in the wild is reported at 8 to 13 years and probably averages 8 to 10; a female of at least 18 years was reported killed by hunters on Vancouver Island. Cougars may live as long as 20 years in captivity. Causes of death in the wild include disability and disease, competition with other cougars, starvation, accidents, and, where allowed, hunting. The feline immunodeficiency virus is well-adapted to the cougar.
Conservation
The cougar has been listed as Least Concern on the IUCN Red List since 2008. However, it is also listed on CITES Appendix II. Hunting it is prohibited in California, Costa Rica, Honduras, Nicaragua, Guatemala, Panama, Venezuela, Colombia, French Guiana, Suriname, Bolivia, Brazil, Chile, Paraguay, Uruguay and most of Argentina. Hunting is regulated in Canada, Mexico, Peru, and the United States. Establishing wildlife corridors and protecting sufficient range areas are critical for the sustainability of cougar populations. Research simulations showed that it faces a low extinction risk in areas larger than . Between one and four new individuals entering a population per decade markedly increases persistence, thus highlighting the importance of habitat corridors.
The Florida panther population is afforded protection under the Endangered Species Act. The Texas Mountain Lion Conservation Project was launched in 2009 and aimed at raising local people's awareness of the status and ecological role of the cougar and mitigating conflict between landowners and cougars.
The cougar is threatened by habitat loss, habitat fragmentation, and depletion of its prey base due to poaching. Hunting is legal in the western United States. In Florida, heavy traffic causes frequent accidents involving cougars. Highways are a major barrier to the dispersal of cougars. The cougar populations in California are becoming fragmented with the increase in human population and infrastructure growth in the state.
Human–wildlife conflict in proximity of of cougar habitat is pronounced in areas with a median human density of and a median livestock population density of . Conflict is generally lower in areas more than away from roads and away from settlements.
Relationships with humans
Attacks on humans
In North America
Due to the expanding human population, cougar ranges increasingly overlap with areas inhabited by humans. Attacks on humans are very rare, as cougar prey recognition is a learned behavior and they do not generally recognize humans as prey. In a 10-year study in New Mexico of wild cougars who were not habituated to humans, the animals did not exhibit threatening behavior to researchers who approached closely (median distance=18.5 m; 61 feet) except in 6% of cases; of those were females with cubs. Attacks on people, livestock, and pets may occur when a puma habituates to humans or is in a condition of severe starvation. Attacks are most frequent during late spring and summer when juvenile cougars leave their mothers and search for new territory.
Between 1890 and 1990 in North America, there were 53 reported, confirmed attacks on humans, resulting in 48 nonfatal injuries and 10 deaths of humans (the total is greater than 53 because some attacks had more than one victim). By 2004, the count had climbed to 88 attacks and 20 deaths.
Within North America, the distribution of attacks is not uniform. The heavily populated state of California saw a dozen attacks from 1986 to 2004 (after just three from 1890 to 1985), including three fatalities. In March 2024, two brothers in California were attacked by a male cougar, with one being fatally wounded; it was the state's first fatal attack in 20 years. Washington state was the site of a fatal attack in 2018, its first since 1924. Lightly populated New Mexico reported an attack in 2008, the first there since 1974.
As with many predators, a cougar may attack if cornered, if a fleeing human stimulates their instinct to chase, or if a person "plays dead". Standing still may cause the cougar to consider a person easy prey. Exaggerating the threat to the animal through intense eye contact, loud shouting, and any other action to appear larger and more menacing, may make the animal retreat. Fighting back with sticks and rocks, or even bare hands, is often effective in persuading an attacking cougar to disengage.
When cougars do attack, they usually employ their characteristic neck bite, attempting to position their teeth between the vertebrae and into the spinal cord. Neck, head, and spinal injuries are common and sometimes fatal. Children are at greatest risk of attack and least likely to survive an encounter. Detailed research into attacks before 1991 showed that 64% of all victims – and almost all fatalities – were children. The same study showed the highest proportion of attacks to have occurred in British Columbia, particularly on Vancouver Island, where cougar populations are especially dense. Preceding attacks on humans, cougars display aberrant behavior, such as activity during daylight hours, a lack of fear of humans, and stalking humans. There have sometimes been incidents of pet cougars mauling people.
Research on new wildlife collars may reduce human-animal conflicts by predicting when and where predatory animals hunt. This may save the lives of humans, pets, and livestock, as well as the lives of these large predatory mammals that are important to the balance of ecosystems.
In South America
Cougars in the southern cone of South America are reputed to be extremely reluctant to attack people; in legend, they defended people against jaguars. The nineteenth-century naturalists Félix de Azara and William Henry Hudson thought that attacks on people, even children or sleeping adults, did not happen. Hudson, citing anecdotal evidence from hunters, claimed that pumas were positively inhibited from attacking people, even in self-defense. Attacks on humans, although exceedingly rare, have occurred.
An early, authenticated, non-fatal case occurred near Lake Viedma, Patagonia, in 1877 when a female mauled the Argentine scientist Francisco P. Moreno; Moreno afterward showed the scars to Theodore Roosevelt. In this instance, however, Moreno had been wearing a guanaco-hide poncho round his neck and head as protection against the cold; in Patagonia the guanaco is the puma's chief prey animal. Another authenticated case occurred in 1997 in Iguazú National Park in northeastern Argentina, when the 20-month-old son of a ranger was killed by a female puma. Forensic analysis found specimens of the child's hair and clothing fibers in the animal's stomach. The coatí is the puma's chief prey in this area. Despite prohibitory signs, coatis are hand-fed by tourists in the park, causing unnatural approximation between cougars and humans. This particular puma had been raised in captivity and released into the wild. On March 13, 2012, Erica Cruz, a 23-year-old shepherdess was found dead in a mountainous area near Rosario de Lerma, Salta Province, in northwestern Argentina. Claw incisions, which severed a jugular vein, indicated that the attacker was a felid; differential diagnosis ruled out other possible perpetrators. There were no bite marks on the victim, who had been herding goats. In 2019 in Córdoba Province, Argentina an elderly man was badly injured by a cougar after he attempted to defend his dog from it, while in neighboring Chile a 28-year-old woman was attacked and killed in Corral, in Los Ríos Region, on October 20, 2020.
Fatal attacks by other carnivores, such as feral dogs, can be misattributed to cougars without appropriate forensic knowledge.
Predation on domestic animals
During the early years of ranching, cougars were considered on par with wolves in destructiveness. According to figures in Texas in 1990, 86 calves (0.0006% of Texas's 13.4 million cattle and calves), 253 mohair goats, 302 mohair kids, 445 sheep (0.02% of Texas's 2 million sheep and lambs) and 562 lambs (0.04% of Texas's 1.2 million lambs) were confirmed to have been killed by cougars that year. In Nevada in 1992, cougars were confirmed to have killed nine calves, one horse, four foals, five goats, 318 sheep, and 400 lambs. In both reports, sheep were the most frequently attacked. Some instances of surplus killing have resulted in the deaths of 20 sheep in one attack. A cougar's killing bite is applied to the back of the neck, head, or throat and the cat inflicts puncture marks with its claws usually seen on the sides and underside of the prey, sometimes also shredding the prey as it holds on. Coyotes also typically bite the throat, but the work of a cougar is generally clean, while bites inflicted by coyotes and dogs leave ragged edges. The size of the tooth puncture marks also helps distinguish kills made by cougars from those made by smaller predators.
Remedial hunting appears to have the paradoxical effect of increased livestock predation and complaints of human-cougar conflicts. In a 2013 study, the most important predictor of cougar problems was the remedial hunting of cougars the previous year. Each additional cougar on the landscape increased predation and human-cougar complaints by 5%, but each animal killed during the previous year increased complaints by 50%. The effect had a dose-response relationship with very heavy (100% removal of adult cougars) remedial hunting, leading to a 150–340% increase in livestock and human conflicts. This effect is attributed to the removal of older cougars that have learned to avoid people and their replacement by younger males that react differently to humans. Remedial hunting enables younger males to enter the former territories of the older animals. Predation by cougars on dogs "is widespread, but occurs at low frequencies".
In mythology
The grace and power of the cougar have been widely admired in the cultures of the indigenous peoples of the Americas. The Inca city of Cusco is reported to have been designed in the shape of a cougar, and the animal also gave its name to both Inca regions and people. The Moche people often represented the cougar in their ceramics. The sky and thunder god of the Inca, Viracocha, has been associated with the animal.
In North America, mythological descriptions of the cougar have appeared in the stories of the Hocąk language ("Ho-Chunk" or "Winnebago") of Wisconsin and Illinois and the Cheyenne, among others. To the Apache and Walapai of the Southwestern United States, the wail of the cougar was a harbinger of death. The Algonquins and Ojibwe believe that the cougar lived in the underworld and was wicked, whereas it was a sacred animal among the Cherokee.
| Biology and health sciences | Carnivora | null |
78214 | https://en.wikipedia.org/wiki/Permaculture | Permaculture | Permaculture is an approach to land management and settlement design that adopts arrangements observed in flourishing natural ecosystems. It includes a set of design principles derived using whole-systems thinking. It applies these principles in fields such as regenerative agriculture, town planning, rewilding, and community resilience. The term was coined in 1978 by Bill Mollison and David Holmgren, who formulated the concept in opposition to modern industrialized methods, instead adopting a more traditional or "natural" approach to agriculture.
Permaculture has been criticised as being poorly defined and unscientific. Critics have pushed for less reliance on anecdote and extrapolation from ecological first principles, in favor of peer-reviewed research to substantiate productivity claims and to clarify methodology. Peter Harper from the Centre for Alternative Technology suggests that most of what passes for permaculture has no relevance to real problems. Defenders of permaculture reply that researchers have concluded it to be a "sustainable alternative to conventional agriculture", that it "strongly" enhances carbon stocks, soil quality, and biodiversity, making it "an effective tool to promote sustainable agriculture, ensure sustainable production patterns, combat climate change and halt and reverse land degradation and biodiversity loss". They further point out that most of permaculture’s most common methods, such as agroforestry, polycultures, and water harvesting features, are also backed by peer-reviewed research.
Background
History
In 1911, Franklin Hiram King wrote Farmers of Forty Centuries: Or Permanent Agriculture in China, Korea and Japan, describing farming practices of East Asia designed for "permanent agriculture". In 1929, Joseph Russell Smith appended King's term as the subtitle for Tree Crops: A Permanent Agriculture, which he wrote in response to widespread deforestation, plow agriculture, and erosion in the eastern mountains and hill regions of the United States. He proposed the planting of tree fruits and nuts as human and animal food crops that could stabilize watersheds and restore soil health. Smith saw the world as an inter-related whole and suggested mixed systems of trees with understory crops. This book inspired individuals such as Toyohiko Kagawa who pioneered forest farming in Japan in the 1930s. Another pioneer, George Washington Carver, advocated for practices now common in permaculture, including the use of crop rotation to restore nitrogen to the soil and repair damaged farmland, in his work at the Tuskegee Institute between 1896 and his death in 1947.
In his 1964 book Water for Every Farm, the Australian agronomist and engineer P. A. Yeomans advanced a definition of permanent agriculture as one that can be sustained indefinitely. Yeomans introduced both an observation-based approach to land use in Australia in the 1940s and in the 1950s the Keyline Design as a way of managing the supply and distribution of water in semi-arid regions. Other early influences include Stewart Brand's works, Ruth Stout and Esther Deans, who pioneered no-dig gardening, and Masanobu Fukuoka who, in the late 1930s in Japan, began advocating no-till orchards and gardens and natural farming.
In the late 1960s, Bill Mollison, senior lecturer in Environmental Psychology at University of Tasmania, and David Holmgren, graduate student at the then Tasmanian College of Advanced Education started developing ideas about stable agricultural systems on the southern Australian island of Tasmania. Their recognition of the unsustainable nature of modern industrialized methods and their inspiration from Tasmanian Aboriginal and other traditional practises were critical to their formulation of permaculture. In their view, industrialized methods were highly dependent on non-renewable resources, and were additionally poisoning land and water, reducing biodiversity, and removing billions of tons of topsoil from previously fertile landscapes. They responded with permaculture. This term was first made public with the publication of their 1978 book Permaculture One.
Following the publication of Permaculture One, Mollison responded to widespread enthusiasm for the work by traveling and teaching a three-week program that became known as the Permaculture Design Course. It addressed the application of permaculture design to growing in major climatic and soil conditions, to the use of renewable energy and natural building methods, and to "invisible structures" of human society. He found ready audiences in Australia, New Zealand, the USA, Britain, and Europe, and from 1985 also reached the Indian subcontinent and southern Africa. By the early 1980s, the concept had broadened from agricultural systems towards sustainable human habitats and at the 1st Intl. Permaculture Convergence, a gathering of graduates of the PDC held in Australia, the curriculum was formalized and its format shortened to two weeks. After Permaculture One, Mollison further refined and developed the ideas while designing hundreds of properties. This led to the 1988 publication of his global reference work, Permaculture: A Designers Manual. Mollison encouraged graduates to become teachers and set up their own institutes and demonstration sites. Critics suggest that this success weakened permaculture's social aspirations of moving away from industrial social forms. They argue that the self-help model (akin to franchising) has had the effect of creating market-focused social relationships that the originators initially opposed.
Foundational ethics
The ethics on which permaculture builds are:
"Care of the Earth: Provision for all life systems to continue and multiply".
"Care of people: Provision for people to access those resources necessary for their existence".
"Setting limits to population and consumption: By governing our own needs, we can set resources aside to further the above principles".
Mollison's 1988 formulation of the third ethic was restated by Holmgren in 2002 as "Set limits to consumption and reproduction, and redistribute surplus" and is elsewhere condensed to "share the surplus".
Permaculture emphasizes patterns of landscape, function, and species assemblies. It determines where these elements should be placed so they can provide maximum benefit to the local environment. Permaculture maximizes synergy of the final design. The focus of permaculture, therefore, is not on individual elements, but rather on the relationships among them. The aim is for the whole to become greater than the sum of its parts, minimizing waste, human labour, and energy input, and to and maximize benefits through synergy.
Permaculture design is founded in replicating or imitating natural patterns found in ecosystems because these solutions have emerged through evolution over thousands of years and have proven to be effective. As a result, the implementation of permaculture design will vary widely depending on the region of the Earth it is located in. Because permaculture's implementation is so localized and place specific, scientific literature for the field is lacking or not always applicable. Design principles derive from the science of systems ecology and the study of pre-industrial examples of sustainable land use.
A core theme of permaculture is the idea of "people care". Seeking prosperity begins within a local community or culture that can apply the tenets of permaculture to sustain an environment that supports them and vice versa. This is in contrast to typical modern industrialized societies, where locality and generational knowledge is often overlooked in the pursuit of wealth or other forms of societal leverage.
Theory
Design principles
Holmgren articulated twelve permaculture design principles in his Permaculture: Principles and Pathways Beyond Sustainability:
Observe and interact: Take time to engage with nature to design solutions that suit a particular situation.
Catch and store energy: Develop systems that collect resources at peak abundance for use in times of need.
Obtain a yield: Emphasize projects that generate meaningful rewards.
Apply self-regulation and accept feedback: Discourage inappropriate activity to ensure that systems function well.
Use and value renewable resources and services: Make the best use of nature's abundance: reduce consumption and dependence on non-renewable resources.
Produce no waste: Value and employ all available resources: waste nothing.
Design from patterns to details: Observe patterns in nature and society and use them to inform designs, later adding details.
Integrate rather than segregate: Proper designs allow relationships to develop between design elements, allowing them to work together to support each other.
Use small and slow solutions: Small and slow systems are easier to maintain, make better use of local resources, and produce more sustainable outcomes.
Use and value diversity: Diversity reduces system-level vulnerability to threats and fully exploits its environment.
Use edges and value the marginal: The border between things is where the most interesting events take place. These are often the system's most valuable, diverse, and productive elements.
Creatively use and respond to change: A positive impact on inevitable change comes from careful observation, followed by well-timed intervention.
Guilds
A guild is a mutually beneficial group of species that form a part of the larger ecosystem. Within a guild each species of insect or plant provides a unique set of diverse services that work in harmony. Plants may be grown for food production, drawing nutrients from deep in the soil through tap roots, balancing nitrogen levels in the soil (legumes), for attracting beneficial insects to the garden, and repelling undesirable insects or pests. There are several types of guilds, such as community function guilds, mutual support guilds, and resource partitioning guilds.
Community function guilds group species based on a specific function or niche that they fill in the garden. Examples of this type of guild include plants that attract a particular beneficial insect or plants that restore nitrogen to the soil. These types of guilds are aimed at solving specific problems which may arise in a garden, such as infestations of harmful insects and poor nutrition in the soil.
Establishment guilds are commonly used when working to establish target species (the primary vegetables, fruits, herbs, etc. you want to be established in your garden) with the support of pioneer species (plants that will help the target species succeed). For example, in temperate climates, plants such as comfrey (as a weed barrier and dynamic accumulator), lupine (as a nitrogen fixer), and daffodil (as a gopher deterrent) can together form a guild for a fruit tree. As the tree matures, the support plants will likely eventually be shaded out and can be used as compost.
Mature guilds form once your target species are established. For example, if the tree layer of your landscape closes its canopy, sun-loving support plants will be shaded out and die. Shade loving medicinal herbs such as ginseng, Black Cohosh, and goldenseal can be planted as an understory.
Mutual support guilds group species together that are complementary by working together and supporting each other. This guild may include a plant that fixes nitrogen, a plant that hosts insects that are predators to pests, and another plant that attracts pollinators.
Resource partitioning guilds group species based on their abilities to share essential resources with one another through a process of niche differentiation. An example of this type of guild includes placing a fibrous- or shallow-rooted plant next to a tap-rooted plant so that they draw from different levels of soil nutrients.
Zones
Zones intelligently organize design elements in a human environment based on the frequency of human use and plant or animal needs. Frequently manipulated or harvested elements of the design are located close to the house in zones 1 and 2. Manipulated elements located further away are used less frequently. Zones are numbered from 0 to 5 based on positioning.
Zone 0
The house, or home center. Here permaculture principles aim to reduce energy and water needs harnessing natural resources such as sunlight, to create a harmonious, sustainable environment in which to live and work. Zone 0 is an informal designation, not specifically defined in Mollison's book.
Zone 1
The zone nearest to the house, the location for those elements in the system that require frequent attention, or that need to be visited often, such as salad crops, herb plants, soft fruit like strawberries or raspberries, greenhouse and cold frames, propagation area, worm compost bin for kitchen waste, etc. Raised beds are often used in Zone 1 in urban areas.
Zone 2
This area is used for siting perennial plants that require less frequent maintenance, such as occasional weed control or pruning, including currant bushes and orchards, pumpkins, sweet potato, etc. Also, a good place for beehives, larger-scale composting bins, etc.
Zone 3
The area where main crops are grown, both for domestic use and for trade purposes. After establishment, care and maintenance required are fairly minimal (provided mulches and similar things are used), such as watering or weed control maybe once a week.
Zone 4
A semi-wild area, mainly used for forage and collecting wild plants as well as production of timber for construction or firewood.
Zone 5
A wilderness area. Humans do not intervene in zone 5 apart from observing natural ecosystems and cycles. This zone hosts a natural reserve of bacteria, molds, and insects that can aid the zones above it.
Edge effect
The edge effect in ecology is the increased diversity that results when two habitats meet. Permaculturists argue that these places can be highly productive. An example of this is a coast. Where land and sea meet is a rich area that meets a disproportionate percentage of human and animal needs. This idea is reflected in permacultural designs by using spirals in herb gardens, or creating ponds that have wavy undulating shorelines rather than a simple circle or oval (thereby increasing the amount of edge for a given area). On the other hand, in a keyhole bed, edges are minimized to avoid wasting space and effort.
Common practices
Hügelkultur
Hügelkultur is the practice of burying wood to increase soil water retention. The porous structure of wood acts like a sponge when decomposing underground. During the rainy season, sufficient buried wood can absorb enough water to sustain crops through the dry season. This technique is a traditional practice that has been developed over centuries in Europe and has been recently adopted by permaculturalists. The Hügelkultur technique can be implemented through building mounds on the ground as well as in raised garden beds. In raised beds, the practice "imitates natural nutrient cycling found in wood decomposition and the high water-holding capacities of organic detritus, while also improving bed structure and drainage properties." This is done by placing wood material (e.g. logs and sticks) in the bottom of the bed before piling organic soil and compost on top. A study comparing the water retention capacities of Hügel raised beds to non-Hügel beds determined that Hügel beds are both lower maintenance and more efficient in the long term by requiring less irrigation.
Sheet mulching
Mulch is a protective cover placed over soil. Mulch material includes leaves, cardboard, and wood chips. These absorb rain, reduce evaporation, provide nutrients, increase soil organic matter, create habitat for soil organisms, suppress weed growth and seed germination, moderate diurnal temperature swings, protect against frost, and reduce erosion. Sheet mulching or lasagna gardening is a gardening technique that attempts to mimic the leaf cover that is found on forest floors.
No-till gardening
Edward Faulkner's 1943 book Plowman's Folly, King's 1946 pamphlet "Is Digging Necessary?", A. Guest's 1948 book "Gardening without Digging", and Fukuoka's "Do Nothing Farming" all advocated forms of no-till or no-dig gardening. No-till gardening seeks to minimise disturbance to the soil community so as to maintain soil structure and organic matter.
Cropping practices
Low-effort permaculture favours perennial crops which do not require tilling and planting every year. Annual crops inevitably require more cultivation. They can be incorporated into permaculture by using traditional techniques such as crop rotation, intercropping, and companion planting so that pests and weeds of individual annual crop species do not build up, and minerals used by specific crop plants do not become successively depleted.
Companion planting aims to make use of beneficial interactions between species of cultivated plants. Such interactions include pest control, pollination, providing habitat for beneficial insects, and maximizing use of space; all of these may help to increase productivity.
Rainwater harvesting
Rainwater harvesting is the accumulation and storage of rainwater for reuse before it runs off or reaches the aquifer. It has been used to provide drinking water, water for livestock, and water for irrigation, as well as other typical uses. Rainwater collected from the roofs of houses and local institutions can make an important contribution to the availability of drinking water. It can supplement the water table and increase urban greenery. Water collected from the ground, sometimes from areas which are specially prepared for this purpose, is called stormwater harvesting.
Greywater is wastewater generated from domestic activities such as laundry, dishwashing, and bathing, which can be recycled for uses such as landscape irrigation and constructed wetlands. Greywater is largely sterile, but not potable (drinkable).
Keyline design is a technique for maximizing the beneficial use of water resources. It was developed in Australia by farmer and engineer P. A. Yeomans. Keyline refers to a contour line extending in both directions from a keypoint. Plowing above and below the keyline provides a watercourse that directs water away from a purely downhill course to reduce erosion and encourage infiltration. It is used in designing drainage systems.
Compost production
Vermicomposting is a common practice in permaculture. The practice involves using earthworms, such as red wigglers, to break down green and brown waste. The worms produce worm castings, which can be used to organically fertilize the garden. Worms are also introduced to garden beds, helping to aerate the soil and improve water retention. Worms may multiply quickly if provided conditions are ideal. For example, a permaculture farm in Cuba began with 9 tiger worms in 2001 and 15 years later had a population of over 500,000. The worm castings are particularly useful as part of a seed starting mix and regular fertilizer. Worm castings are reportedly more successful than conventional compost for seed starting.
Sewage or blackwater contains human or animal waste. It can be composted, producing biogas and manure. Human waste can be sourced from a composting toilet, outhouse or dry bog (rather than a plumbed toilet).
Economising on space
Space can be saved in permaculture gardens with techniques such as herb spirals which group plants closely together. A herb spiral, invented by Mollison, is a round cairn of stones packed with earth at the base and sand higher up; sometimes there is a small pond on the south side (in the northern hemisphere). The result is a series of microclimate zones, wetter at the base, drier at the top, warmer and sunnier on the south side, cooler and drier to the north. Each herb is planted in the zone best suited to it.
Domesticated animals
Domesticated animals are often incorporated into site design. Activities that contribute to the system include: foraging to cycle nutrients, clearing fallen fruit, weed maintenance, spreading seeds, and pest maintenance. Nutrients are cycled by animals, transformed from their less digestible form (such as grass or twigs) into more nutrient-dense manure.
Multiple animals can contribute, including cows, goats, chickens, geese, turkey, rabbits, and worms. An example is chickens who can be used to scratch over the soil, thus breaking down the topsoil and using fecal matter as manure. Factors such as timing and habits are critical. For example, animals require much more daily attention than plants.
Fruit trees
Masanobu Fukuoka experimented with no-pruning methods on his family farm in Japan, finding that trees which were never pruned could grow well, whereas previously-pruned trees often died when allowed to grow without further pruning. He felt that this reflected the Tao-philosophy of Wú wéi, meaning no action against nature or "do-nothing" farming. He claimed yields comparable to intensive arboriculture with pruning and chemical fertilisation.
Applications
Agroforestry
Agroforestry uses the interactive benefits from combining trees and shrubs with crops or livestock. It combines agricultural and forestry technologies to create more diverse, productive, profitable, healthy and sustainable land-use systems. Trees or shrubs are intentionally used within agricultural systems, or non-timber forest products are cultured in forest settings.
Forest gardens
Forest gardens or food forests are permaculture systems designed to mimic natural forests. Forest gardens incorporate processes and relationships that the designers understand to be valuable in natural ecosystems. A mature forest ecosystem is organised into layers with constituents such as trees, understory, ground cover, soil, fungi, insects, and other animals. Because plants grow to different heights, a diverse community of organisms can occupy a relatively small space, each at a different layer.
Rhizosphere: Root layers within the soil. The major components of this layer are the soil and the organisms that live within it such as plant roots and zomes (including root crops such as potatoes and other edible tubers), fungi, insects, nematodes, and earthworms.
Soil surface/groundcover: Overlaps with the herbaceous layer and the groundcover layer; however plants in this layer grow much closer to the ground, densely fill bare patches, and typically can tolerate some foot traffic. Cover crops retain soil and lessen erosion, along with green manures that add nutrients and organic matter, especially nitrogen.
Herbaceous layer: Plants that die back to the ground every winter, if cold enough. No woody stems. Many beneficial plants such as culinary and medicinal herbs are in this layer; whether annuals, biennials, or perennials.
Shrub layer: woody perennials of limited height. Includes most berry bushes.
Understory layer: trees that flourish under the canopy.
The canopy: the tallest trees. Large trees dominate, but typically do not saturate the area, i.e., some patches are devoid of trees.
Vertical layer: climbers or vines, such as runner beans and lima beans (vine varieties).
Suburban and urban permaculture
The fundamental element of suburban and urban permaculture is the efficient utilization of space. Wildfire journal suggests using methods such as the keyhole garden which require little space. Neighbors can collaborate to increase the scale of transformation, using sites such as recreation centers, neighborhood associations, city programs, faith groups, and schools. Columbia, an ecovillage in Portland, Oregon, consisting of 37 apartment condominiums, influenced its neighbors to implement permaculture principles, including in front-yard gardens. Suburban permaculture sites such as one in Eugene, Oregon, include rainwater catchment, edible landscaping, removing paved driveways, turning a garage into living space, and changing a south side patio into passive solar.
Vacant lot farms are community-managed farm sites, but are often seen by authorities as temporary rather than permanent. For example, Los Angeles' South Central Farm (1994–2006), one of the largest urban gardens in the United States, was bulldozed with approval from property owner Ralph Horowitz, despite community protest.
The possibilities and challenges for suburban or urban permaculture vary with the built environment around the world. For example, land is used more ecologically in Jaisalmer, India than in American planned cities such as Los Angeles:
Marine systems
Permaculture derives its origin from agriculture, although the same principles, especially its foundational ethics, can also be applied to mariculture, particularly seaweed farming. In Marine Permaculture, artificial upwelling of cold, deep ocean water is induced. When an attachment substrate is provided in association with such an upwelling, and kelp sporophytes are present, a kelp forest ecosystem can be established (since kelp needs the cool temperatures and abundant dissolved macronutrients present in such an environment). Microalgae proliferate as well. Marine forest habitat is beneficial for many fish species, and the kelp is a renewable resource for food, animal feed, medicines and various other commercial products. It is also a powerful tool for carbon fixation.
The upwelling can be powered by renewable energy on location. Vertical mixing has been reduced due to ocean stratification effects associated with climate change. Reduced vertical mixing and marine heatwaves have decimated seaweed ecosystems in many areas. Marine permaculture mitigates this by restoring some vertical mixing and preserves these important ecosystems. By preserving and regenerating habitat offshore on a platform, marine permaculture employs natural processes to regenerate marine life.
Grazing
Grazing is blamed for much destruction. However, when grazing is modeled after nature, it can have the opposite effect. Cell grazing is a system of grazing in which herds or flocks are regularly and systematically moved to fresh range with the intent to maximize forage quality and quantity. Sepp Holzer and Joel Salatin have shown how grazing can start ecological succession or prepare ground for planting. Allan Savory's holistic management technique has been likened to "a permaculture approach to rangeland management". One variation is conservation grazing, where the primary purpose of the animals is to benefit the environment and the animals are not necessarily used for meat, milk or fiber. Sheep can replace lawn mowers. Goats and sheep can eat invasive plants.
Natural building
Natural building involves using a range of building systems and materials that apply permaculture principles. The focus is on durability and the use of minimally processed, plentiful, or renewable resources, as well as those that, while recycled or salvaged, produce healthy living environments and maintain indoor air quality. For example, cement, a common building material, emits carbon dioxide and is harmful to the environment while natural building works with the environment, using materials that are biodegradable, such as cob, adobe, rammed earth (unburnt clay), and straw bale (which insulates as well as modern synthetic materials).
Issues
Intellectual property
Trademark and copyright disputes surround the word permaculture. Mollison's books claimed on the copyright page, "The contents of this book and the word PERMACULTURE are copyright." Eventually Mollison acknowledged that he was mistaken and that no copyright protection existed.
In 2000, Mollison's U.S.-based Permaculture Institute sought a service mark for the word permaculture when used in educational services such as conducting classes, seminars, or workshops. The service mark would have allowed Mollison and his two institutes to set enforceable guidelines regarding how permaculture could be taught and who could teach it, particularly with relation to the PDC, despite the fact that he had been certifying teachers since 1993. This attempt failed and was abandoned in 2001. Mollison's application for trademarks in Australia for the terms "Permaculture Design Course" and "Permaculture Design" was withdrawn in 2003. In 2009 he sought a trademark for "Permaculture: A Designers' Manual" and "Introduction to Permaculture", the names of two of his books. These applications were withdrawn in 2011. Australia has never authorized a trademark for the word permaculture.
Methodology
Permaculture has been criticised as being poorly defined and unscientific. Critics have pushed for less reliance on anecdote and extrapolation from ecological first principles in favor of peer-reviewed research to substantiate productivity claims and to clarify methodology. Peter Harper from the Centre for Alternative Technology suggests that most of what passes for permaculture is irrelevant to real problems. Harper notes that British organic farmers are "embarrassed or openly derisive" of permaculture, while the permaculture expert Robert Kourik found the supposed advantages of "less- or no-work gardening, bountiful yields, and the soft fuzzy glow of knowing that the garden will ... live on without you" were often illusory. Harper found "many permacultures" are based on ideas ranging from practical farming techniques to "bullshit ... no more than charming cultural graces."
Defenders respond that permaculture is not yet a mainstream scientific tradition and lacks the resources of mainstream industrial agriculture. Rafter Ferguson and Sarah Lovell point out that permaculturalists rarely engage with mainstream research in agroecology, agroforestry, or ecological engineering, and claim that mainstream science has an elitist or pro-corporate bias. Julius Krebs and Sonja Bach argue in Sustainability that there is "scientific evidence for all twelve [of Holmgren's] principles".
In 2017, Ferguson and Lovell presented sociological and demographic data from 36 self-described American permaculture farms. The farms were well diversified, with a median effective number of enterprises per farm of 3.6 (out of a maximum of 6 in the analysis method used). Business strategies included small mixed farms, integrated producers of perennial and animal crops, mixes of production and services, livestock, and service-based businesses. Median household income ($38,750) was less than either the national median household income ($51,017) or the national median farm household income ($68,680).
A 2019 study by Hirschfeld and Van Acker found that adopting permaculture consistently encouraged the cultivation of perennials, crop diversity, landscape heterogeneity, and nature conservation. They discovered that grass-roots adopters were "remarkably consistent" in their implementation of permaculture, leading them to conclude that the movement could exert influence over positive agroecological transitions.
In 2024, Reiff and colleagues stated that permaculture is a "sustainable alternative to conventional agriculture", and that it "strongly" enhances carbon stocks, soil quality, and biodiversity, making it "an effective tool to promote sustainable agriculture, ensure sustainable production patterns, combat climate change and halt and reverse land degradation and biodiversity loss." They point out that most of permaculture's commonest methods, such as agroforestry, polycultures, and water harvesting features, are backed by peer-reviewed research.
| Technology | Agriculture_2 | null |
78227 | https://en.wikipedia.org/wiki/Thalamus | Thalamus | The thalamus (: thalami; from Greek θάλαμος, "chamber") is a large mass of gray matter on the lateral walls of the third ventricle forming the dorsal part of the diencephalon (a division of the forebrain). Nerve fibers project out of the thalamus to the cerebral cortex in all directions, known as the thalamocortical radiations, allowing hub-like exchanges of information. It has several functions, such as the relaying of sensory and motor signals to the cerebral cortex and the regulation of consciousness, sleep, and alertness.
Anatomically, it is a paramedian symmetrical structure of two halves (left and right), within the vertebrate brain, situated between the cerebral cortex and the midbrain. It forms during embryonic development as the main product of the diencephalon, as first recognized by the Swiss embryologist and anatomist Wilhelm His Sr. in 1893.
Anatomy
The thalamus is a paired structure of gray matter about four centimetres long, located in the forebrain which is superior to the midbrain, near the center of the brain with nerve fibers projecting out to the cerebral cortex in all directions. In fact, almost all thalamic neurons (with the notable exception of the thalamic reticular nucleus) project to the cerebral cortex, and every region of the cortex so far studied has been found to innervate the thalamus.
Each of the thalami may be subdivided into at least 30 nuclei, giving a total of at least 60 for the whole thalamus.
Estimates of the volume of the whole thalamus vary. A post-mortem study of 10 people with average age 71 years found average volume 13.68 cm. In a study of 12 healthy males with average age 17 years, MRI scans showed mean whole thalamus volume 8.68cm.
The medial surface of the thalamus constitutes the upper part of the lateral wall of the third ventricle, and is connected to the corresponding surface of the opposite thalamus by a flattened gray band, the interthalamic adhesion.
The lateral part of the thalamus is the neothalamus, the phylogenetically newest part of the thalamus, which includes the lateral nuclei, the pulvinar nuclei and the medial and lateral geniculate nuclei.
The surface of the thalamus is covered by two layers of white matter, the stratum zonale covers the dorsal surface, and the external medullary lamina covers the lateral surface. (This stratum zonale should not be confused with the stratum zonale of the superior colliculus.) Within the thalamus the internal medullary lamina divides the nuclei into anterior, medial, and lateral groups.
Derivatives of the diencephalon include the dorsally-located epithalamus (essentially the habenula and annexes) and the perithalamus (prethalamus) containing the zona incerta and the thalamic reticular nucleus. Due to their different ontogenetic origins, the epithalamus and the perithalamus are formally distinguished from the thalamus proper. The metathalamus is made up of the lateral geniculate and medial geniculate nuclei.
The thalamus comprises a system of lamellae (made up of myelinated fibers) that separate different thalamic subparts. Other areas are defined by distinct clusters of neurons, such as the periventricular nucleus, the intralaminar elements, the "nucleus limitans", and others. These latter structures, different in structure from the major part of the thalamus, have been grouped together into the allothalamus as opposed to the isothalamus. This distinction simplifies the global description of the thalamus.
Thalamic nuclei
The principal subdivision of the thalamus into nucleus groups is the trisection of each thalamus (left and right) by a Y-shaped internal medullary lamina. This trisection divides each thalamus into anterior, medial and lateral groups of nuclei. The medial group is subdivided into the medial dorsal nucleus and midline group. The lateral group is subdivided into ventral, pulvinar, lateral dorsal, lateral posterior and metathalamus. The ventral group is further subdivided into ventral anterior, ventral lateral and ventral posterior.
The interior medullary lamina is subdivided into intralaminar nuclei. Additional structures are the reticular nucleus (which envelops the lateral thalamus), the stratum zonale, and the interthalamic adhesion.
Combining these division principles yields the following hierarchy, which is subject to many further subdivisions.
anterior group
medial group
medial dorsal nucleus
midline group
lateral group
ventral group
ventral anterior group
ventral lateral group
ventral posterior group
pulvinar group
lateral dorsal nucleus
lateral posterior nucleus
metathalamus
lateral geniculate nucleus
medial geniculate nucleus
intralaminar group
reticular nucleus
stratum zonale
interthalamic adhesion
The term "lateral nuclear group" is used with two meanings. It can mean either the complete set of nuclei in the lateral "third" of the trisection by the lamina, or the subset which excludes the ventral group and the geniculate nuclei.
Blood supply
The thalamus derives its blood supply from a number of arteries: the polar artery (posterior communicating artery), paramedian thalamic-subthalamic arteries, inferolateral (thalamogeniculate) arteries, and posterior (medial and lateral) choroidal arteries. These are all branches of the posterior cerebral artery.
Some people have the artery of Percheron, which is a rare anatomic variation in which a single arterial trunk arises from the posterior cerebral artery to supply both parts of the thalamus.
Connections
The thalamus has many connections to the hippocampus via the mammillothalamic tract. This tract comprises the mammillary bodies and fornix.
The thalamus is connected to the cerebral cortex via the thalamocortical radiations.
The spinothalamic tract is a sensory pathway originating in the spinal cord. It transmits information to the thalamus about pain, temperature, itch and crude touch. There are two main parts: the lateral spinothalamic tract, which transmits pain and temperature, and the anterior (or ventral) spinothalamic tract, which transmits crude touch and pressure.
Function
The thalamus has multiple functions, and is generally believed to act as a relay station, or hub, relaying information between different subcortical areas and the cerebral cortex. In particular, every sensory system (with the exception of the olfactory system) includes a thalamic nucleus that receives sensory signals and sends them to the associated primary cortical area.
For the visual system, for example, inputs from the retina are sent to the lateral geniculate nucleus of the thalamus, which in turn projects to the visual cortex in the occipital lobe. Similarly the medial geniculate nucleus acts as a key auditory relay between the inferior colliculus of the midbrain and the primary auditory cortex. The ventral posterior nucleus is a key somatosensory relay, which sends touch and proprioceptive information to the primary somatosensory cortex. In rodents, proprioceptive information of head and whisker movements is integrated already at the thalamic level.
The thalamus is believed to both process sensory information as well as relay it—each of the primary sensory relay areas receives strong feedback connections from the cerebral cortex.
The thalamus also plays an important role in regulating states of sleep, and wakefulness. Thalamic nuclei have strong reciprocal connections with the cerebral cortex, forming thalamo-cortico-thalamic circuits that are believed to be involved with consciousness. The thalamus plays a major role in regulating arousal, the level of awareness, and activity. Damage to the thalamus can lead to permanent coma.
The role of the thalamus in the more anterior pallidal and nigral territories in the basal ganglia system disturbances is recognized but still poorly understood. The contribution of the thalamus to vestibular or to tectal functions is almost ignored. The thalamus has been thought of as a "relay" that simply forwards signals to the cerebral cortex. Newer research suggests that thalamic function is more selective. Many different functions are linked to various regions of the thalamus. This is the case for many of the sensory systems (except for the olfactory system), such as the auditory, somatic, visceral, gustatory and visual systems where localized lesions provoke specific sensory deficits. A major role of the thalamus is support of motor and language systems, and much of the circuitry implicated for these systems is shared.
The thalamus is functionally connected to the hippocampus as part of the extended hippocampal system at the thalamic anterior nuclei. With respect to spatial memory and spatial sensory datum they are crucial for human episodic event memory. The thalamic region's connection to the medial temporal lobe provides differentiation of the functioning of recollective and familiarity memory.
The neuronal information processes necessary for motor control were proposed as a network involving the thalamus as a subcortical motor center. Through investigations of the anatomy of the brains of primates the nature of the interconnected tissues of the cerebellum to the multiple motor cortices suggested that the thalamus fulfills a key function in providing the specific channels from the basal ganglia and cerebellum to the cortical motor areas. In an investigation of the saccade and antisaccade motor response in three monkeys, the thalamic regions were found to be involved in the generation of antisaccade eye-movement (that is, the ability to inhibit the reflexive jerking movement of the eyes in the direction of a presented stimulus).
Recent research suggests that the mediodorsal thalamus (MD) may play a broader role in cognition. Specifically, the mediodorsal thalamus may "amplify the connectivity (signaling strength) of just the circuits in the cortex appropriate for the current context and thereby contribute to the flexibility (of the mammalian brain) to make complex decisions by wiring the many associations on which decisions depend into weakly connected cortical circuits." Researchers found that "enhancing MD activity magnified the ability of mice to "think," driving down by more than 25 percent their error rate in deciding which conflicting sensory stimuli to follow to find the reward."
Development
The thalamic complex is composed of the perithalamus (or prethalamus, previously also known as ventral thalamus), the mid-diencephalic organiser (which forms later the zona limitans intrathalamica (ZLI) ) and the thalamus (dorsal thalamus). The development of the thalamus can be subdivided into three steps.
The thalamus is the largest structure deriving from the embryonic diencephalon, the posterior part of the forebrain situated between the midbrain and the cerebrum.
Early brain development
After neurulation, the early developmental stage (primordium) of the prethalamus and the thalamus is induced within the neural tube. Data from different vertebrate model organisms support a model in which the interaction between two transcription factors, Fez and Otx, is of decisive importance. Fez is expressed in the prethalamus, and functional experiments show that Fez is required for prethalamus formation. Posteriorly, OTX1 and OTX2 abut the expression domain of Fez and are required for proper development of the thalamus.
Formation of progenitor domains
Early in thalamic development two progenitor domains form, a caudal domain, and a rostral domain. The caudal domain gives rise to all of the glutamatergic neurons in the adult thalamus while the rostral domain gives rise to all of the GABAergic neurons in the adult thalamus.
The formation of the mid-diencephalic organiser (MDO)
At the interface between the expression domains of Fez and Otx, the mid-diencephalic organizer (MDO, also called the ZLI organiser) is induced within the thalamic anlage. The MDO is the central signalling organizer in the thalamus. A lack of the organizer leads to the absence of the thalamus. The MDO matures from ventral to dorsal during development. Members of the sonic hedgehog (SHH) family and of the Wnt family are the main principal signals emitted by the MDO.
Besides its importance as signalling center, the organizer matures into the morphological structure of the zona limitans intrathalamica (ZLI).
Maturation and parcellation of the thalamus
After its induction, the MDO starts to orchestrate the development of the thalamic anlage by release of signalling molecules such as SHH. In mice, the function of signaling at the MDO has not been addressed directly due to a complete absence of the diencephalon in SHH mutants.
Studies in chicks have shown that SHH is both necessary and sufficient for thalamic gene induction. In zebrafish, it was shown that the expression of two SHH genes, SHH-a and SHH-b (formerly described as twhh) mark the MDO territory, and that SHH signaling is sufficient for the molecular differentiation of both the prethalamus and the thalamus but is not required for their maintenance and SHH signaling from the MDO/alar plate is sufficient for the maturation of prethalamic and thalamic territory while ventral Shh signals are dispensable.
The exposure to SHH leads to differentiation of thalamic neurons. SHH signaling from the MDO induces a posterior-to-anterior wave of expression the proneural gene Neurogenin1 in the major (caudal) part of the thalamus, and Ascl1 (formerly Mash1) in the remaining narrow stripe of rostral thalamic cells immediately adjacent to the MDO, and in the prethalamus.
This zonation of proneural gene expression leads to the differentiation of glutamatergic relay neurons from the Neurogenin1+ precursors and of GABAergic inhibitory neurons from the Ascl1+ precursors. In fish, selection of these alternative neurotransmitter fates is controlled by the dynamic expression of Her6 the homolog of HES1. Expression of this hairy-like bHLH transcription factor, which represses Neurogenin but is required for Ascl1, is progressively lost from the caudal thalamus but maintained in the prethalamus and in the stripe of rostral thalamic cells. In addition, studies on chick and mice have shown that blocking the Shh pathway leads to absence of the rostral thalamus and substantial decrease of the caudal thalamus. The rostral thalamus will give rise to the reticular nucleus mainly whereby the caudal thalamus will form the relay thalamus and will be further subdivided in the thalamic nuclei.
In humans, a common genetic variation in the promoter region of the serotonin transporter (the SERT-long and -short allele: 5-HTTLPR) has been shown to affect the development of several regions of the thalamus in adults. People who inherit two short alleles (SERT-ss) have more neurons and a larger volume in the pulvinar and possibly the limbic regions of the thalamus. Enlargement of the thalamus provides an anatomical basis for why people who inherit two SERT-ss alleles are more vulnerable to major depression, post-traumatic stress disorder, and suicide.
Clinical significance
A thalamus damaged by a stroke can lead to thalamic pain syndrome, which involves a one-sided burning or aching sensation often accompanied by mood swings. Bilateral ischemia of the area supplied by the paramedian artery can cause serious problems including akinetic mutism, and be accompanied by oculomotor problems. A related concept is thalamocortical dysrhythmia. The occlusion of the artery of Percheron can lead to a bilateral thalamus infarction.
Korsakoff syndrome stems from damage to the mammillary body, the mammillothalamic fasciculus or the thalamus.
Fatal familial insomnia is a hereditary prion disease in which degeneration of the thalamus occurs, causing the patient to gradually lose their ability to sleep and progressing to a state of total insomnia, which invariably leads to death. In contrast, damage to the thalamus can result in coma.
Atrophy of the thalamus is an indicator of the start of multiple sclerosis. Thalamic volume loss by atrophy, is also significantly shown in sporadic frontotemporal dementia, noted in the anterior-dorsal thickness.
Microstimulation of the posterior portion of the ventral medial thalamic nucleus can be used to evoke pain, temperature and visceral sensations.
Additional images
| Biology and health sciences | Nervous system | Biology |
78471 | https://en.wikipedia.org/wiki/Estuary | Estuary | An estuary is a partially enclosed coastal body of brackish water with one or more rivers or streams flowing into it, and with a free connection to the open sea. Estuaries form a transition zone between river environments and maritime environments and are an example of an ecotone. Estuaries are subject both to marine influences such as tides, waves, and the influx of saline water, and to fluvial influences such as flows of freshwater and sediment. The mixing of seawater and freshwater provides high levels of nutrients both in the water column and in sediment, making estuaries among the most productive natural habitats in the world.
Most existing estuaries formed during the Holocene epoch with the flooding of river-eroded or glacially scoured valleys when the sea level began to rise about 10,000–12,000 years ago. Estuaries are typically classified according to their geomorphological features or to water-circulation patterns. They can have many different names, such as bays, harbors, lagoons, inlets, or sounds, although some of these water bodies do not strictly meet the above definition of an estuary and could be fully saline.
Many estuaries suffer degeneration from a variety of factors including soil erosion, deforestation, overgrazing, overfishing and the filling of wetlands. Eutrophication may lead to excessive nutrients from sewage and animal wastes; pollutants including heavy metals, polychlorinated biphenyls, radionuclides and hydrocarbons from sewage inputs; and diking or damming for flood control or water diversion.
Definition
The word "estuary" is derived from the Latin word aestuarium meaning tidal inlet of the sea, which in itself is derived from the term aestus, meaning tide. There have been many definitions proposed to describe an estuary. The most widely accepted definition is: "a semi-enclosed coastal body of water, which has a free connection with the open sea, and within which seawater is measurably diluted with freshwater derived from land drainage". However, this definition excludes a number of coastal water bodies such as coastal lagoons and brackish seas.
A more comprehensive definition of an estuary is "a semi-enclosed body of water connected to the sea as far as the tidal limit or the salt intrusion limit and receiving freshwater runoff; however the freshwater inflow may not be perennial, the connection to the sea may be closed for part of the year and tidal influence may be negligible". This broad definition also includes fjords, lagoons, river mouths, and tidal creeks. An estuary is a dynamic ecosystem having a connection to the open sea through which the sea water enters with the rhythm of the tides. The effects of tides on estuaries can show nonlinear effects on the movement of water which can have important impacts on the ecosystem and waterflow. The seawater entering the estuary is diluted by the fresh water flowing from rivers and streams. The pattern of dilution varies between different estuaries and depends on the volume of freshwater, the tidal range, and the extent of evaporation of the water in the estuary.
Classification based on geomorphology
Drowned river valleys
Drowned river valleys are also known as coastal plain estuaries. In places where the sea level is rising relative to the land, sea water progressively penetrates into river valleys and the topography of the estuary remains similar to that of a river valley. This is the most common type of estuary in temperate climates. Well-studied estuaries include the Severn Estuary in the United Kingdom and the Ems Dollard along the Dutch-German border.
The width-to-depth ratio of these estuaries is typically large, appearing wedge-shaped (in cross-section) in the inner part and broadening and deepening seaward. Water depths rarely exceed . Examples of this type of estuary in the U.S. are the Hudson River, Chesapeake Bay, and Delaware Bay along the Mid-Atlantic coast, and Galveston Bay and Tampa Bay along the Gulf Coast.
Lagoon-type or bar-built
Bar-built estuaries are found in a place where the deposition of sediment has kept pace with rising sea levels so that the estuaries are shallow and separated from the sea by sand spits or barrier islands. They are relatively common in tropical and subtropical locations.
These estuaries are semi-isolated from ocean waters by barrier beaches (barrier islands and barrier spits). Formation of barrier beaches partially encloses the estuary, with only narrow inlets allowing contact with the ocean waters. Bar-built estuaries typically develop on gently sloping plains located along tectonically stable edges of continents and marginal sea coasts. They are extensive along the Atlantic and Gulf coasts of the U.S. in areas with active coastal deposition of sediments and where tidal ranges are less than . The barrier beaches that enclose bar-built estuaries have been developed in several ways:
building up of offshore bars by wave action, in which sand from the seafloor is deposited in elongated bars parallel to the shoreline,
reworking of sediment discharge from rivers by a wave, current, and wind action into beaches, overwash flats, and dunes,
engulfment of mainland beach ridges (ridges developed from the erosion of coastal plain sediments around 5000 years ago) due to sea level rise and resulting in the breaching of the ridges and flooding of the coastal lowlands, forming shallow lagoons,
elongation of barrier spits from the erosion of headlands due to the action of longshore currents, with the spits growing in the direction of the littoral drift.
Fjord-type
Fjords were formed where Pleistocene glaciers deepened and widened existing river valleys so that they become U-shaped in cross-sections. At their mouths there are typically rocks, bars or sills of glacial deposits, which have the effects of modifying the estuarine circulation.
Fjord-type estuaries are formed in deeply eroded valleys formed by glaciers. These U-shaped estuaries typically have steep sides, rock bottoms, and underwater sills contoured by glacial movement. The estuary is shallowest at its mouth, where terminal glacial moraines or rock bars form sills that restrict water flow. In the upper reaches of the estuary, the depth can exceed . The width-to-depth ratio is generally small. In estuaries with very shallow sills, tidal oscillations only affect the water down to the depth of the sill, and the waters deeper than that may remain stagnant for a very long time, so there is only an occasional exchange of the deep water of the estuary with the ocean. If the sill depth is deep, water circulation is less restricted, and there is a slow but steady exchange of water between the estuary and the ocean. Fjord-type estuaries can be found along the coasts of Alaska, the Puget Sound region of western Washington state, British Columbia, eastern Canada, Greenland, Iceland, New Zealand, Chile, and Norway.
Tectonically produced
These estuaries are formed by subsidence or land cut off from the ocean by land movement associated with faulting, volcanoes, and landslides. Inundation from eustatic sea-level rise during the Holocene Epoch has also contributed to the formation of these estuaries. There are only a small number of tectonically produced estuaries; one example is the San Francisco Bay, which was formed by the crustal movements of the San Andreas Fault system causing the inundation of the lower reaches of the Sacramento and San Joaquin rivers.
Classification based on water circulation
Salt wedge
In this type of estuary, river output greatly exceeds marine input and tidal effects have minor importance. Freshwater floats on top of the seawater in a layer that gradually thins as it moves seaward. The denser seawater moves landward along the bottom of the estuary, forming a wedge-shaped layer that is thinner as it approaches land. As a velocity difference develops between the two layers, shear forces generate internal waves at the interface, mixing the seawater upward with the freshwater. An examples of a salt wedge estuary is Mississippi River and the Mandovi estuary in Goa during the monsoon period.
Partially mixed
As tidal forcing increases, river output becomes less than the marine input. Here, current induced turbulence causes mixing of the whole water column such that salinity varies more longitudinally rather than vertically, leading to a moderately stratified condition. Examples include the Chesapeake Bay and Narragansett Bay.
Well-mixed
Tidal mixing forces exceed river output, resulting in a well-mixed water column and the disappearance of the vertical salinity gradient. The freshwater-seawater boundary is eliminated due to the intense turbulent mixing and eddy effects. The lower reaches of Delaware Bay and the Raritan River in New Jersey are examples of vertically homogeneous estuaries.
Inverse
Inverse estuaries occur in dry climates where evaporation greatly exceeds the inflow of freshwater. A salinity maximum zone is formed, and both riverine and oceanic water flow close to the surface towards this zone. This water is pushed downward and spreads along the bottom in both the seaward and landward direction. Examples of an inverse estuary are Spencer Gulf, South Australia, Saloum River and Casamance River, Senegal.
Intermittent
Estuary type varies dramatically depending on freshwater input, and is capable of changing from a wholly marine embayment to any of the other estuary types.
Physiochemical variation
The most important variable characteristics of estuary water are the concentration of dissolved oxygen, salinity and sediment load. There is extreme spatial variability in salinity, with a range of near-zero at the tidal limit of tributary rivers to 3.4% at the estuary mouth. At any one point, the salinity will vary considerably over time and seasons, making it a harsh environment for organisms. Sediment often settles in intertidal mudflats which are extremely difficult to colonize. No points of attachment exist for algae, so vegetation based habitat is not established. Sediment can also clog feeding and respiratory structures of species, and special adaptations exist within mudflat species to cope with this problem. Lastly, dissolved oxygen variation can cause problems for life forms. Nutrient-rich sediment from human-made sources can promote primary production life cycles, perhaps leading to eventual decay removing the dissolved oxygen from the water; thus hypoxic or anoxic zones can develop.
Implications of eutrophication on estuaries
Effects of eutrophication on biogeochemical cycles
Nitrogen is often the lead cause of eutrophication in estuaries in temperate zones. During a eutrophication event, biogeochemical feedback decreases the amount of available silica. These feedbacks also increase the supply of nitrogen and phosphorus, creating conditions where harmful algal blooms can persist. Given the now off-balance nitrogen cycle, estuaries can be driven to phosphorus limitation instead of nitrogen limitation. Estuaries can be severely impacted by an unbalanced phosphorus cycle, as phosphorus interacts with nitrogen and silica availability.
With an abundance of nutrients in the ecosystem, plants and algae overgrow and eventually decompose, which produce a significant amount of carbon dioxide. While releasing into the water and atmosphere, these organisms are also intaking all or nearly all of the available oxygen creating a hypoxic environment and unbalanced oxygen cycle. The excess carbon in the form of can lead to low pH levels and ocean acidification, which is more harmful for vulnerable coastal regions like estuaries.
Effects of eutrophication on estuarine plants
Eutrophication has been seen to negatively impact many plant communities in estuarine ecosystems. Salt marshes are a type of ecosystem in some estuaries that have been negatively impacted by eutrophication. Cordgrass vegetation dominates the salt marsh landscape. Excess nutrients allow the plants to grow at greater rates in above ground biomass, however less energy is allocated to the roots since nutrients is abundant. This leads to a lower biomass in the vegetation below ground which destabilizes the banks of the marsh causing increased rates of erosion. A similar phenomenon occurs in mangrove swamps, which are another potential ecosystem in estuaries. An increase in nitrogen causes an increase in shoot growth and a decrease in root growth. Weaker root systems cause a mangrove tree to be less resilient in seasons of drought, which can lead to the death of the mangrove. This shift in above ground and below ground biomass caused by eutrophication could hindered plant success in these ecosystems.
Effects of eutrophication on estuarine animals
Across all biomes, eutrophication often results in plant death but the impacts do not end there. Plant death alters the entire food web structure which can result in the death of animals within the afflicted biome. Estuaries are hotspots for biodiversity, containing a majority of commercial fish catch, making the impacts of eutrophication that much greater within estuaries. Some specific estuarine animals feel the effects of eutrophication more strongly than others. One example is the whitefish species from the European Alps. Eutrophication reduced the oxygen levels in their habitats so greatly that whitefish eggs could not survive, causing local extinctions. However, some animals, such as carnivorous fish, tend to do well in nutrient-enriched environments and can benefit from eutrophication. This can be seen in populations of bass or pikes.
Effects of eutrophication on human activities
Eutrophication can affect many marine habitats which can lead to economic consequences. The commercial fishing industry relies upon estuaries for approximately 68 percent of their catch by value because of the great biodiversity of this ecosystem. During an algal bloom, fishermen have noticed a significant increase in the quantity of fish. A sudden increase in primary productivity causes spikes in fish populations which leads to more oxygen being utilized. It is the continued deoxygenation of the water that then causes a decline in fish populations. These effects can begin in estuaries and have a wide effect on the surrounding water bodies. In turn, this can decrease fishing industry sales in one area and across the country. Production in 2016 from recreational and commercial fishing contributes billions of dollars to the United States' gross domestic product (GDP). A decrease in production within this industry can affect any of the 1.7 million people the fishing industry employs yearly across the United States.
Implications for marine life
Estuaries are incredibly dynamic systems, where temperature, salinity, turbidity, depth and flow all change daily in response to the tides. This dynamism makes estuaries highly productive habitats, but also make it difficult for many species to survive year-round. As a result, estuaries large and small experience strong seasonal variation in their fish communities. In winter, the fish community is dominated by hardy marine residents, and in summer a variety of marine and anadromous fishes move into and out of estuaries, capitalizing on their high productivity. Estuaries provide a critical habitat to a variety of species that rely on estuaries for life-cycle completion. Pacific Herring (Clupea pallasii) are known to lay their eggs in estuaries and bays, surfperch give birth in estuaries, juvenile flatfish and rockfish migrate to estuaries to rear, and anadromous salmonids and lampreys use estuaries as migration corridors. Also, migratory bird populations, such as the black-tailed godwit, rely on estuaries.
Two of the main challenges of estuarine life are the variability in salinity and sedimentation. Many species of fish and invertebrates have various methods to control or conform to the shifts in salt concentrations and are termed osmoconformers and osmoregulators. Many animals also burrow to avoid predation and to live in a more stable sedimental environment. However, large numbers of bacteria are found within the sediment which has a very high oxygen demand. This reduces the levels of oxygen within the sediment often resulting in partially anoxic conditions, which can be further exacerbated by limited water flow.
Phytoplankton are key primary producers in estuaries. They move with the water bodies and can be flushed in and out with the tides. Their productivity is largely dependent upon the turbidity of the water. The main phytoplankton present are diatoms and dinoflagellates which are abundant in the sediment.
A primary source of food for many organisms on estuaries, including bacteria, is detritus from the settlement of the sedimentation.
Human impact
Of the thirty-two largest cities in the world in the early 1990s, twenty-two were located on estuaries.
As ecosystems, estuaries are under threat from human activities such as pollution and overfishing. They are also threatened by sewage, coastal settlement, land clearance and much more. Estuaries are affected by events far upstream, and concentrate materials such as pollutants and sediments. Land run-off and industrial, agricultural, and domestic waste enter rivers and are discharged into estuaries. Contaminants can be introduced which do not disintegrate rapidly in the marine environment, such as plastics, pesticides, furans, dioxins, phenols and heavy metals.
Such toxins can accumulate in the tissues of many species of aquatic life in a process called bioaccumulation. They also accumulate in benthic environments, such as estuaries and bay muds: a geological record of human activities of the last century. The elemental composition of biofilm reflect areas of the estuary impacted by human activities, and over time may shift the basic composition of the ecosystem, and the reversible or irreversible changes in the abiotic and biotic parts of the systems from the bottom up.
For example, Chinese and Russian industrial pollution, such as phenols and heavy metals, has devastated fish stocks in the Amur River and damaged its estuary soil.
Estuaries tend to be naturally eutrophic because land runoff discharges nutrients into estuaries. With human activities, land run-off also now includes the many chemicals used as fertilizers in agriculture as well as waste from livestock and humans. Excess oxygen-depleting chemicals in the water can lead to hypoxia and the creation of dead zones. This can result in reductions in water quality, fish, and other animal populations.
Overfishing also occurs. Chesapeake Bay once had a flourishing oyster population that has been almost wiped out by overfishing. Oysters filter these pollutants, and either eat them or shape them into small packets that are deposited on the bottom where they are harmless. Historically the oysters filtered the estuary's entire water volume of excess nutrients every three or four days. Today that process takes almost a year, and sediment, nutrients, and algae can cause problems in local waters.
Some major rivers that run through deserts historically had vast, expansive estuaries that have been reduced to a fraction of their former size, because of dams and diversions. One example is the Colorado River Delta in Mexico, historically covered with marshlands and forests, but now essentially a salt flat.
Examples
Africa
Congo River Estuary
Estuário do Espírito Santo
Gambia River Estuary
Gabon Estuary
Lake St Lucia Estuary
Orange River Estuary
Pungwe River Estuary
Asia
Amur River Estuary
Adyar River Estuary
Dawei River Estuary
Gulf of Ob Estuary
Hangzhou Bay
Hàn River Estuary
Kraburi River Estuary
Meghna River Estuary
Naf River Estuary
Narmada river estuary
Puerto Princesa Underground River
Waeru River Estuary of Chanthaburi Province
Yangtze River estuary
Yenisei Gulf Estuary
Europe
Dee Estuary
Dnieper-Bug Estuary
Exe Estuary
Firth of Clyde
Firth of Forth
Gironde estuary
Golden Horn
Humber
Oder Estuary
Severn Estuary
Shannon Estuary
Solway Firth
Southampton Water
Tagus Estuary
Thames Estuary
The Wash
Unterelbe
Western Scheldt
North America
Albemarle Sound including Outer Banks of North Carolina
Chesapeake Bay including Hampton Roads
Columbia River Estuary
Coos Bay
Delaware Bay
Drake's Estero
East River
Estuary of Saint Lawrence
Fraser River
Galveston Bay
Great Bay
Indian River Lagoon
Laguna de Términos
Laguna Madre
Lake Borgne
Lake Merritt
Long Island Sound
Miramichi Bay
Mississippi River Delta
Lake Pontchartrain
Mobile Bay
Narragansett Bay
Newport Back Bay
New York-New Jersey Harbor
Pamlico Sound including the Outer Banks of North Carolina
Puget Sound
San Francisco Bay
Sarasota Bay
Tampa Bay
Oceania
Avon Heathcote Estuary (Christchurch, New Zealand)
Gippsland Lakes
Port Jackson (Sydney Harbour)
Spencer Gulf
South America
Amazon River
Iguape-Cananéia-Paranaguá estuary lagoon complex
Lagoa dos Patos and Lagoon Mirim
Mearim River
São Marcos Bay
São José Bay
Rio de la Plata
| Physical sciences | Hydrology | null |
1105937 | https://en.wikipedia.org/wiki/Electrical%20wiring | Electrical wiring | Electrical wiring is an electrical installation of cabling and associated devices such as switches, distribution boards, sockets, and light fittings in a structure.
Wiring is subject to safety standards for design and installation. Allowable wire and cable types and sizes are specified according to the circuit operating voltage and electric current capability, with further restrictions on the environmental conditions, such as ambient temperature range, moisture levels, and exposure to sunlight and chemicals.
Associated circuit protection, control, and distribution devices within a building's wiring system are subject to voltage, current, and functional specifications. Wiring safety codes vary by locality, country, or region. The International Electrotechnical Commission (IEC) is attempting to harmonise wiring standards among member countries, but significant variations in design and installation requirements still exist.
Wiring methods
Materials for wiring interior electrical systems in buildings vary depending on:
Intended use and amount of power demand on the circuit
Type of occupancy and size of the building
National and local regulations
Environment in which the wiring must operate.
Wiring systems in a single family home or duplex, for example, are simple, with relatively low power requirements, infrequent changes to the building structure and layout, usually with dry, moderate temperature and non-corrosive environmental conditions. In a light commercial environment, more frequent wiring changes can be expected, large apparatus may be installed and special conditions of heat or moisture may apply. Heavy industries have more demanding wiring requirements, such as very large currents and higher voltages, frequent changes of equipment layout, corrosive, or wet or explosive atmospheres. In facilities that handle flammable gases or liquids, special rules may govern the installation and wiring of electrical equipment in hazardous areas.
Wires and cables are rated by the circuit voltage, temperature rating and environmental conditions (moisture, sunlight, oil, chemicals) in which they can be used. A wire or cable has a voltage (to neutral) rating and a maximum conductor surface temperature rating. The amount of current a cable or wire can safely carry depends on the installation conditions.
The international standard wire sizes are given in the IEC 60228 standard of the International Electrotechnical Commission. In North America, the American Wire Gauge standard for wire sizes is used.
Cables
Modern wiring materials
Modern non-metallic sheathed cables, such as (US and Canadian) Types NMB and NMC, consist of two to four wires covered with thermoplastic insulation, plus a wire for Protective Earthing/Grounding (bonding), surrounded by a flexible plastic jacket. In North America and the UK this conductor is usually bare wire but in the UK it is required that this bare Protective Earth (PE) conductor be sheathed in Green/Yellow insulating tubing where the Cable Sheathing has been removed. Most other jurisdictions now require the Protective Earth conductor to be insulated to the same standard as the current carrying conductors with Green/Yellow insulation.
With some cables the individual conductors are wrapped in paper before the plastic jacket is applied.
Special versions of non-metallic sheathed cables, such as US Type UF, are designed for direct underground burial (often with separate mechanical protection) or exterior use where exposure to ultraviolet radiation (UV) is a possibility. These cables differ in having a moisture-resistant construction, lacking paper or other absorbent fillers, and being formulated for UV resistance.
Rubber-like synthetic polymer insulation is used in industrial cables and power cables installed underground because of its superior moisture resistance.
Insulated cables are rated by their allowable operating voltage and their maximum operating temperature at the conductor surface. A cable may carry multiple usage ratings for applications, for example, one rating for dry installations and another when exposed to moisture or oil.
Generally, single conductor building wire in small sizes is solid wire, since the wiring is not required to be very flexible. Building wire conductors larger than 10 AWG (or about 5 mm2) are stranded for flexibility during installation, but are not sufficiently pliable to use as appliance cord.
Cables for industrial, commercial and apartment buildings may contain many insulated conductors in an overall jacket, with helical tape steel or aluminium armour, or steel wire armour, and perhaps as well an overall PVC or lead jacket for protection from moisture and physical damage. Cables intended for very flexible service or in marine applications may be protected by woven bronze wires. Power or communications cables (e.g., computer networking) that are routed in or through air-handling spaces (plenums) of office buildings are required under the model building code to be either encased in metal conduit, or rated for low flame and smoke production.
For some industrial uses in steel mills and similar hot environments, no organic material gives satisfactory service. Cables insulated with compressed mica flakes are sometimes used. Another form of high-temperature cable is mineral-insulated cable, with individual conductors placed within a copper tube and the space filled with magnesium oxide powder. The whole assembly is drawn down to smaller sizes, thereby compressing the powder. Such cables have a certified fire resistance rating and are more costly than non–fire-rated cable. They have little flexibility and behave more like rigid conduit rather than flexible cables.
The environment of the installed wires determine how much current a cable is permitted to carry. Because multiple conductors bundled in a cable cannot dissipate heat as easily as single insulated conductors, those circuits are always rated at a lower ampacity. Tables in electrical safety codes give the maximum allowable current based on size of conductor, voltage potential, insulation type and thickness, and the temperature rating of the cable itself. The allowable current will also be different for wet or dry locations, for hot (attic) or cool (underground) locations. In a run of cable through several areas, the part with the lowest rating becomes the rating of the overall run.
Cables usually are secured with special fittings where they enter electrical apparatus; this may be a simple screw clamp for jacketed cables in a dry location, or a polymer-gasketed cable connector that mechanically engages the armour of an armoured cable and provides a water-resistant connection. Special cable fittings may be applied to prevent explosive gases from flowing in the interior of jacketed cables, where the cable passes through areas where flammable gases are present. To prevent loosening of the connections of individual conductors of a cable, cables must be supported near their entrance to devices and at regular intervals along their runs. In tall buildings, special designs are required to support the conductors of vertical runs of cable. Generally, only one cable per fitting is permitted, unless the fitting is rated or listed for multiple cables.
Special cable constructions and termination techniques are required for cables installed in ships. Such assemblies are subjected to environmental and mechanical extremes. Therefore, in addition to electrical and fire safety concerns, such cables may also be required to be pressure-resistant where they penetrate a vessel's bulkheads. They must also resist corrosion caused by salt water or salt spray, which is accomplished through the use of thicker, specially constructed jackets, and by tinning the individual wire stands.
In North American practice, for residential and light commercial buildings fed with a single-phase split 120/240 service, an overhead cable from a transformer on a power pole is run to the service entrance point. The cable is a three conductor twisted "triplex" cable with a bare neutral and two insulated conductors, with no overall cable jacket. The neutral conductor is often a supporting "messenger" steel wire, which is used to support the insulated line conductors.
Copper conductors
Electrical devices often use copper conductors because of their properties, including their high electrical conductivity, tensile strength, ductility, creep resistance, corrosion resistance, thermal conductivity, coefficient of thermal expansion, solderability, resistance to electrical overloads, compatibility with electrical insulators, and ease of installation. Copper is used in many types of electrical wiring.
Aluminium conductors
Aluminium wire was common in North American residential wiring from the late 1960s to mid-1970s due to the rising cost of copper. Because of its greater resistivity, aluminium wiring requires larger conductors than copper. For instance, instead of 14 AWG (American wire gauge) copper wire, aluminium wiring would need to be 12 AWG on a typical 15 ampere lighting circuit, though local building codes vary.
Solid aluminium conductors were originally made in the 1960s from a utility-grade aluminium alloy that had undesirable properties for a building wire, and were used with wiring devices intended for copper conductors. These practices were found to cause defective connections and fire hazards. In the early 1970s new aluminium wire made from one of several special alloys was introduced, and all devices – breakers, switches, receptacles, splice connectors, wire nuts, etc. — were specially designed for the purpose. These newer aluminium wires and special designs address problems with junctions between dissimilar metals, oxidation on metal surfaces, and mechanical effects that occur as different metals expand at different rates with increases in temperature.
Unlike copper, aluminium has a tendency to creep or cold-flow under pressure, so older plain steel screw clamped connections could become loose over time. Newer electrical devices designed for aluminium conductors have features intended to compensate for this effect. Unlike copper, aluminium forms an insulating oxide layer on the surface. This is sometimes addressed by coating aluminium conductors with an antioxidant paste (containing zinc dust in a low-residue polybutene base) at joints, or by applying a mechanical termination designed to break through the oxide layer during installation.
Some terminations on wiring devices designed only for copper wire would overheat under heavy current load and cause fires when used with aluminium conductors. Revised standards for wire materials and wiring devices (such as the CO/ALR "copper-aluminium-revised" designation) were developed to reduce these problems. While larger sizes are still used to feed power to electrical panels and large devices, aluminium wiring for residential use has acquired a poor reputation and has fallen out of favour.
Aluminium conductors are still heavily used for bulk power transmission, power distribution, and large feeder circuits with heavy current loads, due to the various advantages they offer over copper wiring. Aluminium conductors both cost and weigh less than copper conductors, so a much larger cross sectional area can be used for the same weight and price. This can compensate for the higher resistance and lower mechanical strength of aluminium, meaning the larger cross sectional area is needed to achieve comparable current capacity and other features. Aluminium conductors must be installed with compatible connectors and special care must be taken to ensure the contact surface does not oxidise.
Raceways and cable runs
Insulated wires may be run in one of several forms between electrical devices. This may be a specialised bendable pipe, called a conduit, or one of several varieties of metal (rigid steel or aluminium) or non-metallic (PVC or HDPE) tubing. Rectangular cross-section metal or PVC wire troughs (North America) or trunking (UK) may be used if many circuits are required. Wires run underground may be run in plastic tubing encased in concrete, but metal elbows may be used in severe pulls. Wiring in exposed areas, for example factory floors, may be run in cable trays or rectangular raceways having lids.
Where wiring, or raceways that hold the wiring, must traverse fire-resistance rated walls and floors, the openings are required by local building codes to be firestopped. In cases where safety-critical wiring must be kept operational during an accidental fire, fireproofing must be applied to maintain circuit integrity in a manner to comply with a product's certification listing. The nature and thickness of any passive fire protection materials used in conjunction with wiring and raceways has a quantifiable impact upon the ampacity derating, because the thermal insulation properties needed for fire resistance also inhibit air cooling of power conductors.
Cable trays are used in industrial areas where many insulated cables are run together. Individual cables can exit the tray at any point, simplifying the wiring installation and reducing the labour cost for installing new cables. Power cables may have fittings in the tray to maintain clearance between the conductors, but small control wiring is often installed without any intentional spacing between cables.
Local electrical regulations may restrict or place special requirements on mixing of voltage levels within one cable tray. Good design practices may segregate, for example, low level measurement or signal cables from trays carrying high power branch circuits, to prevent induction of noise into sensitive circuits.
Since wires run in conduits or underground cannot dissipate heat as easily as in open air, and since adjacent circuits contribute induced currents, wiring regulations give rules to establish the current capacity (ampacity).
Special sealed fittings are used for wiring routed through potentially explosive atmospheres.
Bus bars, bus duct, cable bus
For very high currents in electrical apparatus, and for high currents distributed through a building, bus bars can be used. (The term "bus" is a contraction of the Latin omnibus – meaning "for all".) Each live ("hot") conductor of such a system is a rigid piece of copper or aluminium, usually in flat bars (but sometimes as tubing or other shapes). Open bus bars are never used in publicly accessible areas, although they are used in manufacturing plants and power company switch yards to gain the benefit of air cooling. A variation is to use heavy cables, especially where it is desirable to transpose or "roll" phases.
In industrial applications, conductor bars are often pre-assembled with insulators in grounded enclosures. This assembly, known as bus duct or busway, can be used for connections to large switchgear or for bringing the main power feed into a building. A form of bus duct known as "plug-in bus" is used to distribute power down the length of a building; it is constructed to allow tap-off switches or motor controllers to be installed at designated places along the bus. The big advantage of this scheme is the ability to remove or add a branch circuit without removing voltage from the whole duct.
Bus ducts may have all phase conductors in the same enclosure (non-isolated bus), or may have each conductor separated by a grounded barrier from the adjacent phases (segregated bus). For conducting large currents between devices, a cable bus is used.
For very large currents in generating stations or substations, where it is difficult to provide circuit protection, an isolated-phase bus is used. Each phase of the circuit is run in a separate grounded metal enclosure. The only fault possible is a phase-to-ground fault, since the enclosures are separated. This type of bus can be rated up to 50,000 amperes and up to hundreds of kilovolts (during normal service, not just for faults), but is not used for building wiring in the conventional sense.
Electrical panels
Electrical panels are easily accessible junction boxes used to reroute and switch electrical services. The term is often used to refer to circuit breaker panels or fuseboxes. Local codes can specify physical clearance around the panels.
Degradation by pests
Squirrels, rats, and other rodents may gnaw on unprotected wiring, causing fire and shock hazards. This is especially true of PVC-insulated telephone and computer network cables. Several techniques have been developed to deter these pests, including insulation loaded with pepper dust.
Early wiring methods
The first interior power wiring systems used conductors that were bare or covered with cloth, which were secured by staples to the framing of the building or on running boards. Where conductors went through walls, they were protected with cloth tape. Splices were done similarly to telegraph connections, and soldered for security. Underground conductors were insulated with wrappings of cloth tape soaked in pitch, and laid in wooden troughs which were then buried. Such wiring systems were unsatisfactory because of the danger of electrocution and fire, plus the high labour cost for such installations.
The first electrical codes arose in the 1880s with the commercial introduction of electrical power; however, many conflicting standards existed for the selection of wire sizes and other design rules for electrical installations, and a need was seen to introduce uniformity on the grounds of safety.
Knob and tube (US)
The earliest standardized method of wiring in buildings, in common use in North America from about 1880 to the 1930s, was knob and tube (K&T) wiring: single conductors were run through cavities between the structural members in walls and ceilings, with ceramic tubes forming protective channels through joists and ceramic knobs attached to the structural members to provide air between the wire and the lumber and to support the wires. Since air was free to circulate over the wires, smaller conductors could be used than required in cables. By arranging wires on opposite sides of building structural members, some protection was afforded against short-circuits that can be caused by driving a nail into both conductors simultaneously.
By the 1940s, the labor cost of installing two conductors rather than one cable resulted in a decline in new knob-and-tube installations. However, the US code still allows new K&T wiring installations in special situations (some rural and industrial applications).
Metal-sheathed wires
In the United Kingdom, an early form of insulated cable, introduced in 1896, consisted of two impregnated-paper-insulated conductors in an overall lead sheath. Joints were soldered, and special fittings were used for lamp holders and switches. These cables were similar to underground telegraph and telephone cables of the time. Paper-insulated cables proved unsuitable for interior wiring installations because very careful workmanship was required on the lead sheaths to ensure moisture did not affect the insulation.
A system later invented in the UK in 1908 employed vulcanised-rubber insulated wire enclosed in a strip metal sheath. The metal sheath was bonded to each metal wiring device to ensure earthing continuity.
A system developed in Germany called "Kuhlo wire" used one, two, or three rubber-insulated wires in a brass or lead-coated iron sheet tube, with a crimped seam. The enclosure could also be used as a return conductor. Kuhlo wire could be run exposed on surfaces and painted, or embedded in plaster. Special outlet and junction boxes were made for lamps and switches, made either of porcelain or sheet steel. The crimped seam was not considered as watertight as the Stannos wire used in England, which had a soldered sheath.
A somewhat similar system called "concentric wiring" was introduced in the United States around 1905. In this system, an insulated electrical wire was wrapped with copper tape which was then soldered, forming the grounded (return) conductor of the wiring system. The bare metal sheath, at earth potential, was considered safe to touch. While companies such as General Electric manufactured fittings for the system and a few buildings were wired with it, it was never adopted into the US National Electrical Code. Drawbacks of the system were that special fittings were required, and that any defect in the connection of the sheath would result in the sheath becoming energised.
Other historical wiring methods
Armored cables with two rubber-insulated conductors in a flexible metal sheath were used as early as 1906, and were considered at the time a better method than open knob-and-tube wiring, although much more expensive.
The first rubber-insulated cables for US building wiring were introduced in 1922 with . These were two or more solid copper electrical wires with rubber insulation, plus woven cotton cloth over each conductor for protection of the insulation, with an overall woven jacket, usually impregnated with tar as a protection from moisture. Waxed paper was used as a filler and separator.
Over time, rubber-insulated cables become brittle because of exposure to atmospheric oxygen, so they must be handled with care and are usually replaced during renovations. When switches, socket outlets or light fixtures are replaced, the mere act of tightening connections may cause hardened insulation to flake off the conductors. Rubber insulation further inside the cable often is in better condition than the insulation exposed at connections, due to reduced exposure to oxygen.
The sulfur in vulcanized rubber insulation attacked bare copper wire so the conductors were tinned to prevent this. The conductors reverted to being bare when rubber ceased to be used.
About 1950, PVC insulation and jackets were introduced, especially for residential wiring. About the same time, single conductors with a thinner PVC insulation and a thin nylon jacket (e.g. US Type THN, THHN, etc.) became common.
The simplest form of cable has two insulated conductors twisted together to form a unit. Such non-jacketed cables with two (or more) conductors are used only for extra-low voltage signal and control applications such as doorbell wiring.
Other methods of securing wiring that are now obsolete include:
Re-use of existing gas pipes when converting gas lighting installations to electric lighting. Insulated conductors were pulled through the pipes that had formerly supplied the gas lamps. Although used occasionally, this method risked insulation damage from sharp edges inside the pipe at each joint.
Wood mouldings with grooves cut for single conductor wires, covered by a wooden cap strip. These were prohibited in North American electrical codes by 1928. Wooden moulding was also used to some degree in the UK, but was never permitted by German and Austrian rules.
A system of flexible twin cords supported by glass or porcelain buttons was used near the turn of the 20th century in Europe, but was soon replaced by other methods.
During the first years of the 20th century, various patented forms of wiring system such as Bergman and Peschel tubing were used to protect wiring; these used very thin fibre tubes, or metal tubes which were also used as return conductors.
In Austria, wires were concealed by embedding a rubber tube in a groove in the wall, plastering over it, then removing the tube and pulling wires through the cavity.
Metal moulding systems, with a flattened oval section consisting of a base strip and a snap-on cap channel, were more costly than open wiring or wooden moulding, but could be easily run on wall surfaces. Similar surface mounted raceway wiring systems are still available today.
| Technology | Other components | null |
1106531 | https://en.wikipedia.org/wiki/Intrusive%20rock | Intrusive rock | Intrusive rock is formed when magma penetrates existing rock, crystallizes, and solidifies underground to form intrusions, such as batholiths, dikes, sills, laccoliths, and volcanic necks.
Intrusion is one of the two ways igneous rock can form. The other is extrusion, such as a volcanic eruption or similar event. An intrusion is any body of intrusive igneous rock, formed from magma that cools and solidifies within the crust of the planet. In contrast, an extrusion consists of extrusive rock, formed above the surface of the crust.
Some geologists use the term plutonic rock synonymously with intrusive rock, but other geologists subdivide intrusive rock, by crystal size, into coarse-grained plutonic rock (typically formed deeper in the Earth's crust in batholiths or stocks) and medium-grained subvolcanic or hypabyssal rock (typically formed higher in the crust in dikes and sills).
Classification
Because the solid country rock into which magma intrudes is an excellent insulator, cooling of the magma is extremely slow, and intrusive igneous rock is coarse-grained (phaneritic). However, the rate of cooling is greatest for intrusions at relatively shallow depth, and the rock in such intrusions is often much less coarse-grained than intrusive rock formed at greater depth. Coarse-grained intrusive igneous rocks that form at depth within the Earth are called abyssal or plutonic while those that form near the surface are called subvolcanic or hypabyssal.
Plutonic rocks are classified separately from extrusive igneous rocks, generally on the basis of their mineral content. The relative amounts of quartz, alkali feldspar, plagioclase, and feldspathoid are particularly important in classifying intrusive igneous rocks, and most plutonic rocks are classified by where they fall in the QAPF diagram. Dioritic and gabbroic rocks are further distinguished by whether the plagioclase they contain is sodium-rich, and sodium-poor gabbros are classified by their relative contents of various iron- or magnesium-rich minerals (mafic minerals) such as olivine, hornblende, clinopyroxene, and orthopyroxene, which are the most common mafic minerals in intrusive rock. Rare ultramafic rocks, which contain more than 90% mafic minerals, and carbonatite rocks, containing over 50% carbonate minerals, have their own special classifications.
Hypabyssal rocks resemble volcanic rocks more than they resemble plutonic rocks, being nearly as fine-grained, and are usually assigned volcanic rock names. However, dikes of basaltic composition often show grain sizes intermediate between plutonic and volcanic rock, and are classified as diabases or dolerites. Rare ultramafic hypabyssal rocks called lamprophyres have their own classification scheme.
Characteristics
Intrusive rocks are characterized by large crystal sizes, and as the individual crystals are visible, the rock is called phaneritic. There are few indications of flow in intrusive rocks, since their texture and structure mostly develops in the final stages of crystallization, when flow has ended. Contained gases cannot escape through the overlying strata, and these gases sometimes form cavities, often lined with large, well-shaped crystals. These are particularly common in granites and their presence is described as miarolitic texture. Because their crystals are of roughly equal size, intrusive rocks are said to be equigranular.
Plutonic rocks are less likely than volcanic rocks to show a pronounced porphyritic texture, in which a first generation of large well-shaped crystals are embedded in a fine-grained ground-mass. The minerals of each have formed in a definite order, and each has had a period of crystallization that may be very distinct or may have coincided with or overlapped the period of formation of some of the other ingredients. Earlier crystals originated at a time when most of the rock was still liquid and are more or less perfect. Later crystals are less regular in shape because they were compelled to occupy the spaces left between the already-formed crystals. The former case is said to be idiomorphic (or automorphic); the latter is xenomorphic.
There are also many other characteristics that serve to distinguish plutonic from volcanic rock. For example, the alkali feldspar in plutonic rocks is typically orthoclase, while the higher-temperature polymorph, sanidine, is more common in volcanic rock. The same distinction holds for nepheline varieties. Leucite is common in lavas but very rare in plutonic rocks. Muscovite is confined to intrusions. These differences show the influence of the physical conditions under which crystallization takes place.
Hypabyssal rocks show structures intermediate between those of extrusive and plutonic rocks. They are very commonly porphyritic, vitreous, and sometimes even vesicular. In fact, many of them are petrologically indistinguishable from lavas of similar composition.
Occurrences
Plutonic rocks form 7% of the Earth's current land surface. Intrusions vary widely, from mountain-range-sized batholiths to thin veinlike fracture fillings of aplite or pegmatite.
Batholith: a large irregular discordant intrusion
Chonolith: an irregularly-shaped intrusion with a demonstrable base
Cupola: a dome-shaped projection from the top of a large subterranean intrusion
Dike: a relatively narrow tabular discordant body, often nearly vertical
Laccolith: concordant body with roughly flat base and convex top, usually with a feeder pipe below
Lopolith: concordant body with roughly flat top and a shallow convex base, may have a feeder dike or pipe below
Phacolith: a concordant lens-shaped pluton that typically occupies the crest of an anticline or trough of a syncline
Volcanic pipe or volcanic neck: tubular, roughly vertical body that may have been a feeder vent for a volcano
Sill: a relatively thin tabular concordant body intruded along bedding planes
Stock: a smaller irregular discordant intrusive
Boss: a small stock
| Physical sciences | Igneous rocks | Earth science |
1106830 | https://en.wikipedia.org/wiki/Cell%20culture | Cell culture | Cell culture or tissue culture is the process by which cells are grown under controlled conditions, generally outside of their natural environment. After cells of interest have been isolated from living tissue, they can subsequently be maintained under carefully controlled conditions. They need to be kept at body temperature (37 °C) in an incubator. These conditions vary for each cell type, but generally consist of a suitable vessel with a substrate or rich medium that supplies the essential nutrients (amino acids, carbohydrates, vitamins, minerals), growth factors, hormones, and gases (CO2, O2), and regulates the physio-chemical environment (pH buffer, osmotic pressure, temperature). Most cells require a surface or an artificial substrate to form an adherent culture as a monolayer (one single-cell thick), whereas others can be grown free floating in a medium as a suspension culture. This is typically facilitated via use of a liquid, semi-solid, or solid growth medium, such as broth or agar. Tissue culture commonly refers to the culture of animal cells and tissues, with the more specific term plant tissue culture being used for plants. The lifespan of most cells is genetically determined, but some cell-culturing cells have been 'transformed' into immortal cells which will reproduce indefinitely if the optimal conditions are provided.
In practice, the term "cell culture" now refers to the culturing of cells derived from multicellular eukaryotes, especially animal cells, in contrast with other types of culture that also grow cells, such as plant tissue culture, fungal culture, and microbiological culture (of microbes). The historical development and methods of cell culture are closely interrelated with those of tissue culture and organ culture. Viral culture is also related, with cells as hosts for the viruses.
The laboratory technique of maintaining live cell lines (a population of cells descended from a single cell and containing the same genetic makeup) separated from their original tissue source became more robust in the middle 20th century.
History
The 19th-century English physiologist Sydney Ringer developed salt solutions containing the chlorides of sodium, potassium, calcium and magnesium suitable for maintaining the beating of an isolated animal heart outside the body. In 1885 Wilhelm Roux removed a section of the medullary plate of an embryonic chicken and maintained it in a warm saline solution for several days, establishing the basic principle of tissue culture. In 1907 the zoologist Ross Granville Harrison demonstrated the growth of frog embryonic cells that would give rise to nerve cells in a medium of clotted lymph. In 1913, E. Steinhardt, C. Israeli, and R. A. Lambert grew vaccinia virus in fragments of guinea pig corneal tissue. In 1996, the first use of regenerative tissue was used to replace a small length of urethra, which led to the understanding that the technique of obtaining samples of tissue, growing it outside the body without a scaffold, and reapplying it, can be used for only small distances of less than 1 cm. Ross Granville Harrison, working at Johns Hopkins Medical School and then at Yale University, published results of his experiments from 1907 to 1910, establishing the methodology of tissue culture.
Gottlieb Haberlandt first pointed out the possibilities of the culture of isolated tissues, plant tissue culture. He suggested that the potentialities of individual cells via tissue culture as well as that the reciprocal influences of tissues on one another could be determined by this method. Since Haberlandt's original assertions, methods for tissue and cell culture have been realized, leading to significant discoveries in biology and medicine. He presented his original idea of totipotentiality in 1902, stating that "Theoretically all plant cells are able to give rise to a complete plant." The term tissue culture was coined by American pathologist Montrose Thomas Burrows.
Cell culture techniques were advanced significantly in the 1940s and 1950s to support research in virology. Growing viruses in cell cultures allowed preparation of purified viruses for the manufacture of vaccines. The injectable polio vaccine developed by Jonas Salk was one of the first products mass-produced using cell culture techniques. This vaccine was made possible by the cell culture research of John Franklin Enders, Thomas Huckle Weller, and Frederick Chapman Robbins, who were awarded a Nobel Prize for their discovery of a method of growing the virus in monkey kidney cell cultures. Cell culture has contributed to the development of vaccines for many diseases.
Modern usage
In modern usage, "tissue culture" generally refers to the growth of cells from a tissue from a multicellular organism in vitro. These cells may be cells isolated from a donor organism (primary cells) or an immortalised cell line. The cells are bathed in a culture medium, which contains essential nutrients and energy sources necessary for the cells' survival. Thus, in its broader sense, "tissue culture" is often used interchangeably with "cell culture". On the other hand, the strict meaning of "tissue culture" refers to the culturing of tissue pieces, i.e. explant culture.
Tissue culture is an important tool for the study of the biology of cells from multicellular organisms. It provides an in vitro model of the tissue in a well defined environment which can be easily manipulated and analysed. In animal tissue culture, cells may be grown as two-dimensional monolayers (conventional culture) or within fibrous scaffolds or gels to attain more naturalistic three-dimensional tissue-like structures (3D culture). A 1988 NIH SBIR grant report showed that electrospinning could be used to produce nano- and submicron-scale polymeric fibrous scaffolds specifically intended for use as in vitro cell and tissue substrates. This early use of electrospun fibrous lattices for cell culture and tissue engineering showed that various cell types would adhere to and proliferate upon polycarbonate fibers. It was noted that as opposed to the flattened morphology typically seen in 2D culture, cells grown on the electrospun fibers exhibited a more rounded 3-dimensional morphology generally observed of tissues in vivo.
Plant tissue culture in particular is concerned with the growing of entire plants from small pieces of plant tissue, cultured in medium.
Concepts in mammalian cell culture
Isolation of cells
Cells can be isolated from tissues for ex vivo culture in several ways. Cells can be easily purified from blood; however, only the white cells are capable of growth in culture. Cells can be isolated from solid tissues by digesting the extracellular matrix using enzymes such as collagenase, trypsin, or pronase, before agitating the tissue to release the cells into suspension. Alternatively, pieces of tissue can be placed in growth media, and the cells that grow out are available for culture. This method is known as explant culture.
Cells that are cultured directly from a subject are known as primary cells. With the exception of some derived from tumors, most primary cell cultures have limited lifespan.
An established or immortalized cell line has acquired the ability to proliferate indefinitely either through random mutation or deliberate modification, such as artificial expression of the telomerase gene.
Numerous cell lines are well established as representative of particular cell types.
Maintaining cells in culture
For the majority of isolated primary cells, they undergo the process of senescence and stop dividing after a certain number of population doublings while generally retaining their viability (described as the Hayflick limit).
Aside from temperature and gas mixture, the most commonly varied factor in culture systems is the cell growth medium. Recipes for growth media can vary in pH, glucose concentration, growth factors, and the presence of other nutrients. The growth factors used to supplement media are often derived from the serum of animal blood, such as fetal bovine serum (FBS), bovine calf serum, equine serum, and porcine serum. One complication of these blood-derived ingredients is the potential for contamination of the culture with viruses or prions, particularly in medical biotechnology applications. Current practice is to minimize or eliminate the use of these ingredients wherever possible and use human platelet lysate (hPL). This eliminates the worry of cross-species contamination when using FBS with human cells. hPL has emerged as a safe and reliable alternative as a direct replacement for FBS or other animal serum. In addition, chemically defined media can be used to eliminate any serum trace (human or animal), but this cannot always be accomplished with different cell types. Alternative strategies involve sourcing the animal blood from countries with minimum BSE/TSE risk, such as The United States, Australia and New Zealand, and using purified nutrient concentrates derived from serum in place of whole animal serum for cell culture.
Plating density (number of cells per volume of culture medium) plays a critical role for some cell types. For example, a lower plating density makes granulosa cells exhibit estrogen production, while a higher plating density makes them appear as progesterone-producing theca lutein cells.
Cells can be grown either in suspension or adherent cultures. Some cells naturally live in suspension, without being attached to a surface, such as cells that exist in the bloodstream. There are also cell lines that have been modified to be able to survive in suspension cultures so they can be grown to a higher density than adherent conditions would allow. Adherent cells require a surface, such as tissue culture plastic or microcarrier, which may be coated with extracellular matrix (such as collagen and laminin) components to increase adhesion properties and provide other signals needed for growth and differentiation. Most cells derived from solid tissues are adherent. Another type of adherent culture is organotypic culture, which involves growing cells in a three-dimensional (3-D) environment as opposed to two-dimensional culture dishes. This 3D culture system is biochemically and physiologically more similar to in vivo tissue, but is technically challenging to maintain because of many factors (e.g. diffusion).
Cell culture basal media
There are different kinds of cell culture media which being used routinely in life science including the following:
MEM
DMEM
RPMI 1640
Ham's f-12
IMDM
Leibovitz L-15
DMEM/F-12
GMEM
Components of cell culture media
Typical Growth conditions
Cell line cross-contamination
Cell line cross-contamination can be a problem for scientists working with cultured cells. Studies suggest anywhere from 15 to 20% of the time, cells used in experiments have been misidentified or contaminated with another cell line. Problems with cell line cross-contamination have even been detected in lines from the NCI-60 panel, which are used routinely for drug-screening studies. Major cell line repositories, including the American Type Culture Collection (ATCC), the European Collection of Cell Cultures (ECACC) and the German Collection of Microorganisms and Cell Cultures (DSMZ), have received cell line submissions from researchers that were misidentified by them. Such contamination poses a problem for the quality of research produced using cell culture lines, and the major repositories are now authenticating all cell line submissions. ATCC uses short tandem repeat (STR) DNA fingerprinting to authenticate its cell lines.
To address this problem of cell line cross-contamination, researchers are encouraged to authenticate their cell lines at an early passage to establish the identity of the cell line. Authentication should be repeated before freezing cell line stocks, every two months during active culturing and before any publication of research data generated using the cell lines. Many methods are used to identify cell lines, including isoenzyme analysis, human lymphocyte antigen (HLA) typing, chromosomal analysis, karyotyping, morphology and STR analysis.
One significant cell-line cross contaminant is the immortal HeLa cell line. HeLa contamination was first noted in the early 1960s in non-human culture in the USA. Intraspecies contamination was discovered in nineteen cell lines in the seventies. In 1974, five human cell lines from the Soviet Union were found to be HeLa. A follow-up study analysing 50-odd cell lines indicated that half had HeLa markers, but contaminant HeLa had hybridised with the original cell lines. HeLa cell contamination from air droplets has been reported. HeLa was even unknowingly injected into human subjects by Jonas Salk in a 1978 vaccine trial.
Other technical issues
As cells generally continue to divide in culture, they generally grow to fill the available area or volume. This can generate several issues:
Nutrient depletion in the growth media
Changes in pH of the growth media
Accumulation of apoptotic/necrotic (dead) cells
Cell-to-cell contact can stimulate cell cycle arrest, causing cells to stop dividing, known as contact inhibition.
Cell-to-cell contact can stimulate cellular differentiation.
Genetic and epigenetic alterations, with a natural selection of the altered cells potentially leading to overgrowth of abnormal, culture-adapted cells with decreased differentiation and increased proliferative capacity.
The choice of culture medium might affect the physiological relevance of findings from cell culture experiments due to the differences in the nutrient composition and concentrations. A systematic bias in generated datasets was recently shown for CRISPR and RNAi gene silencing screens, and for metabolic profiling of cancer cell lines. Using a growth medium that better represents the physiological levels of nutrients can improve the physiological relevance of in vitro studies and recently such media types, as Plasmax and Human Plasma Like Medium (HPLM), were developed.
Manipulation of cultured cells
Among the common manipulations carried out on culture cells are media changes, passaging cells, and transfecting cells.
These are generally performed using tissue culture methods that rely on aseptic technique. Aseptic technique aims to avoid contamination with bacteria, yeast, or other cell lines. Manipulations are typically carried out in a biosafety cabinet or laminar flow cabinet to exclude contaminating micro-organisms. Antibiotics (e.g. penicillin and streptomycin) and antifungals (e.g.amphotericin B and Antibiotic-Antimycotic solution) can also be added to the growth media.
As cells undergo metabolic processes, acid is produced and the pH decreases. Often, a pH indicator is added to the medium to measure nutrient depletion.
Media changes
In the case of adherent cultures, the media can be removed directly by aspiration, and then is replaced. Media changes in non-adherent cultures involve centrifuging the culture and resuspending the cells in fresh media.
Passaging cells
Passaging (also known as subculture or splitting cells) involves transferring a small number of cells into a new vessel. Cells can be cultured for a longer time if they are split regularly, as it avoids the senescence associated with prolonged high cell density. Suspension cultures are easily passaged with a small amount of culture containing a few cells diluted in a larger volume of fresh media. For adherent cultures, cells first need to be detached; this is commonly done with a mixture of trypsin-EDTA; however, other enzyme mixes are now available for this purpose. A small number of detached cells can then be used to seed a new culture. Some cell cultures, such as RAW cells are mechanically scraped from the surface of their vessel with rubber scrapers.
Transfection and transduction
Another common method for manipulating cells involves the introduction of foreign DNA by transfection. This is often performed to cause cells to express a gene of interest. More recently, the transfection of RNAi constructs have been realized as a convenient mechanism for suppressing the expression of a particular gene/protein. DNA can also be inserted into cells using viruses, in methods referred to as transduction, infection or transformation. Viruses, as parasitic agents, are well suited to introducing DNA into cells, as this is a part of their normal course of reproduction.
Established human cell lines
Cell lines that originate with humans have been somewhat controversial in bioethics, as they may outlive their parent organism and later be used in the discovery of lucrative medical treatments. In the pioneering decision in this area, the Supreme Court of California held in Moore v. Regents of the University of California that human patients have no property rights in cell lines derived from organs removed with their consent.
It is possible to fuse normal cells with an immortalised cell line. This method is used to produce monoclonal antibodies. In brief, lymphocytes isolated from the spleen (or possibly blood) of an immunised animal are combined with an immortal myeloma cell line (B cell lineage) to produce a hybridoma which has the antibody specificity of the primary lymphocyte and the immortality of the myeloma. Selective growth medium (HA or HAT) is used to select against unfused myeloma cells; primary lymphoctyes die quickly in culture and only the fused cells survive. These are screened for production of the required antibody, generally in pools to start with and then after single cloning.
Cell strains
A cell strain is derived either from a primary culture or a cell line by the selection or cloning of cells having specific properties or characteristics which must be defined. Cell strains are cells that have been adapted to culture but, unlike cell lines, have a finite division potential. Non-immortalized cells stop dividing after 40 to 60 population doublings and, after this, they lose their ability to proliferate (a genetically determined event known as senescence).
Applications of cell culture
Mass culture of animal cell lines is fundamental to the manufacture of viral vaccines and other products of biotechnology. Culture of human stem cells is used to expand the number of cells and differentiate the cells into various somatic cell types for transplantation. Stem cell culture is also used to harvest the molecules and exosomes that the stem cells release for the purposes of therapeutic development.
Biological products produced by recombinant DNA (rDNA) technology in animal cell cultures include enzymes, synthetic hormones, immunobiologicals (monoclonal antibodies, interleukins, lymphokines), and anticancer agents. Although many simpler proteins can be produced using rDNA in bacterial cultures, more complex proteins that are glycosylated (carbohydrate-modified) currently must be made in animal cells. Mammalian cells ensure expressed proteins are folded correctly and possess human-like glycosylation and post-translational modifications. An important example of such a complex protein is the hormone erythropoietin. The cost of growing mammalian cell cultures is high, so research is underway to produce such complex proteins in insect cells or in higher plants, use of single embryonic cell and somatic embryos as a source for direct gene transfer via particle bombardment, transit gene expression and confocal microscopy observation is one of its applications. It also offers to confirm single cell origin of somatic embryos and the asymmetry of the first cell division, which starts the process.
Cell culture is also a key technique for cellular agriculture, which aims to provide both new products and new ways of producing existing agricultural products like milk, (cultured) meat, fragrances, and rhino horn from cells and microorganisms. It is therefore considered one means of achieving animal-free agriculture. It is also a central tool for teaching cell biology.
Cell culture in two dimensions
Research in tissue engineering, stem cells and molecular biology primarily involves cultures of cells on flat plastic dishes. This technique is known as two-dimensional (2D) cell culture, and was first developed by Wilhelm Roux who, in 1885, removed a portion of the medullary plate of an embryonic chicken and maintained it in warm saline for several days on a flat glass plate. From the advance of polymer technology arose today's standard plastic dish for 2D cell culture, commonly known as the Petri dish. Julius Richard Petri, a German bacteriologist, is generally credited with this invention while working as an assistant to Robert Koch. Various researchers today also utilize culturing laboratory flasks, conicals, and even disposable bags like those used in single-use bioreactors.
Aside from Petri dishes, scientists have long been growing cells within biologically derived matrices such as collagen or fibrin, and more recently, on synthetic hydrogels such as polyacrylamide or PEG. They do this in order to elicit phenotypes that are not expressed on conventionally rigid substrates. There is growing interest in controlling matrix stiffness, a concept that has led to discoveries in fields such as:
Stem cell self-renewal
Lineage specification
Cancer cell phenotype
Fibrosis
Hepatocyte function
Mechanosensing
Cell culture in three dimensions
Cell culture in three dimensions has been touted as "Biology's New Dimension". At present, the practice of cell culture remains based on varying combinations of single or multiple cell structures in 2D. Currently, there is an increase in use of 3D cell cultures in research areas including drug discovery, cancer biology, regenerative medicine, nanomaterials assessment and basic life science research. 3D cell cultures can be grown using a scaffold or matrix, or in a scaffold-free manner. Scaffold based cultures utilize an acellular 3D matrix or a liquid matrix. Scaffold-free methods are normally generated in suspensions. There are a variety of platforms used to facilitate the growth of three-dimensional cellular structures including scaffold systems such as hydrogel matrices and solid scaffolds, and scaffold-free systems such as low-adhesion plates, nanoparticle facilitated magnetic levitation, hanging drop plates, and rotary cell culture. Culturing cells in 3D leads to wide variation in gene expression signatures and partly mimics tissues in the physiological states. A 3D cell culture model showed cell growth similar to that of in vivo than did a monolayer culture, and all three cultures were capable of sustaining cell growth. As 3D culturing has been developed it turns out to have a great potential to design tumors models and investigate malignant transformation and metastasis, 3D cultures can provide aggerate tool for understanding changes, interactions, and cellular signaling.
3D cell culture in scaffolds
Eric Simon, in a 1988 NIH SBIR grant report, showed that electrospinning could be used to produce nano- and submicron-scale polystyrene and polycarbonate fibrous scaffolds specifically intended for use as in vitro cell substrates. This early use of electrospun fibrous lattices for cell culture and tissue engineering showed that various cell types including Human Foreskin Fibroblasts (HFF), transformed Human Carcinoma (HEp-2), and Mink Lung Epithelium (MLE) would adhere to and proliferate upon polycarbonate fibers. It was noted that, as opposed to the flattened morphology typically seen in 2D culture, cells grown on the electrospun fibers exhibited a more histotypic rounded 3-dimensional morphology generally observed in vivo.
3D cell culture in hydrogels
As the natural extracellular matrix (ECM) is important in the survival, proliferation, differentiation and migration of cells, different hydrogel culture matrices mimicking natural ECM structure are seen as potential approaches to in vivo–like cell culturing. Hydrogels are composed of interconnected pores with high water retention, which enables efficient transport of substances such as nutrients and gases. Several different types of hydrogels from natural and synthetic materials are available for 3D cell culture, including animal ECM extract hydrogels, protein hydrogels, peptide hydrogels, polymer hydrogels, and wood-based nanocellulose hydrogel.
3D Cell Culturing by Magnetic Levitation
The 3D Cell Culturing by Magnetic Levitation method (MLM) is the application of growing 3D tissue by inducing cells treated with magnetic nanoparticle assemblies in spatially varying magnetic fields using neodymium magnetic drivers and promoting cell to cell interactions by levitating the cells up to the air/liquid interface of a standard petri dish. The magnetic nanoparticle assemblies consist of magnetic iron oxide nanoparticles, gold nanoparticles, and the polymer polylysine. 3D cell culturing is scalable, with the capability for culturing 500 cells to millions of cells or from single dish to high-throughput low volume systems.
Tissue culture and engineering
Cell culture is a fundamental component of tissue culture and tissue engineering, as it establishes the basics of growing and maintaining cells in vitro.
The major application of human cell culture is in stem cell industry, where mesenchymal stem cells can be cultured and cryopreserved for future use. Tissue engineering potentially offers dramatic improvements in low cost medical care for hundreds of thousands of patients annually.
Vaccines
Vaccines for polio, measles, mumps, rubella, and chickenpox are currently made in cell cultures. Due to the H5N1 pandemic threat, research into using cell culture for influenza vaccines is being funded by the United States government. Novel ideas in the field include recombinant DNA-based vaccines, such as one made using human adenovirus (a common cold virus) as a vector,
and novel adjuvants.
Cell co-culture
The technique of co-culturing is used to study cell crosstalk between two or more types of cells on a plate or in a 3D matrix. The cultivation of different stem cells and the interaction of immune cells can be investigated in an in vitro model similar to biological tissue. Since most tissues contain more than one type of cell, it is important to evaluate their interaction in a 3D culture environment to gain a better understanding of their interaction and to introduce mimetic tissues. There are two types of co-culturing: direct and indirect. While direct interaction involves cells being in direct contact with each other in the same culture media or matrix, indirect interaction involves different environments, allowing signaling and soluble factors to participate.
Cell differentiation in tissue models during interaction between cells can be studied using the Co-Cultured System to simulate cancer tumors, to assess the effect of drugs on therapeutic trials, and to study the effect of drugs on therapeutic trials. The co-culture system in 3D models can predict the response to chemotherapy and endocrine therapy if the microenvironment defines biological tissue for the cells.
A co-culture method is used in tissue engineering to generate tissue formation with multiple cells interacting directly.
Cell culture in microfluidic device
Microfluidics technique is developed systems that can perform a process in a flow which are usually in a scale of micron. Microfluidics chip are also known as Lab-on-a-chip and they are able to have continuous procedure and reaction steps with spare amount of reactants and space. Such systems enable the identification and isolation of individual cells and molecules when combined with appropriate biological assays and high-sensitivity detection techniques.
Organ-on-a-chip
OoC systems mimic and control the microenvironment of the cells by growing tissues in microfluidics. Combining tissue engineering, biomaterials fabrication, and cell biology, it offers the possibility of establishing a biomimetic model for studying human diseases in the laboratory. In recent years, 3D cell culture science has made significant progress, leading to the development of OoC. OoC is considered as a preclinical step that benefits pharmaceutical studies, drug development and disease modeling. OoC is an important technology that can bridge the gap between animal testing and clinical studies and also by the advances that the science has achieved could be a replace for in vivo studies for drug delivery and pathophysiological studies.
Culture of non-mammalian cells
Besides the culture of well-established immortalised cell lines, cells from primary explants of a plethora of organisms can be cultured for a limited period of time before senescence occurs (see Hayflick's limit). Cultured primary cells have been extensively used in research, as is the case of fish keratocytes in cell migration studies.
Plant cell culture methods
Plant cell cultures are typically grown as cell suspension cultures in a liquid medium or as callus cultures on a solid medium. The culturing of undifferentiated plant cells and calli requires the proper balance of the plant growth hormones auxin and cytokinin.
Insect cell culture
Cells derived from Drosophila melanogaster (most prominently, Schneider 2 cells) can be used for experiments which may be hard to do on live flies or larvae, such as biochemical studies or studies using siRNA. Cell lines derived from the army worm Spodoptera frugiperda, including Sf9 and Sf21, and from the cabbage looper Trichoplusia ni, High Five cells, are commonly used for expression of recombinant proteins using baculovirus.
Bacterial and yeast culture methods
For bacteria and yeasts, small quantities of cells are usually grown on a solid support that contains nutrients embedded in it, usually a gel such as agar, while large-scale cultures are grown with the cells suspended in a nutrient broth.
Viral culture methods
The culture of viruses requires the culture of cells of mammalian, plant, fungal or bacterial origin as hosts for the growth and replication of the virus. Whole wild type viruses, recombinant viruses or viral products may be generated in cell types other than their natural hosts under the right conditions. Depending on the species of the virus, infection and viral replication may result in host cell lysis and formation of a viral plaque.
Common cell lines
Human cell lines
DU145 (prostate cancer)
H295R (adrenocortical cancer)
HeLa (cervical cancer)
KBM-7 (chronic myelogenous leukemia)
LNCaP (prostate cancer)
MCF-7 (breast cancer)
MDA-MB-468 (breast cancer)
PC3 (prostate cancer)
SaOS-2 (bone cancer)
SH-SY5Y (neuroblastoma, cloned from a myeloma)
T-47D (breast cancer)
THP-1 (acute myeloid leukemia)
U-87 MG (glioblastoma)
National Cancer Institute's 60 cancer cell line panel (NCI60)
Animal cell lines
Vero (African green monkey Chlorocebus kidney epithelial cell line)
BHK21 cell (Baby Hambster Kidney)
MDBK cell (Madin-Darby Bovine Kidney)
DF-1cell (chicken fibroblast)
Mouse cell lines
MC3T3 (embryonic calvarium)
Rat tumor cell lines
GH3 (pituitary tumor)
PC12 (pheochromocytoma)
Plant cell lines
Tobacco BY-2 cells (kept as cell suspension culture, they are model system of plant cell)
Other species cell lines
Dog MDCK kidney epithelial
Xenopus A6 kidney epithelial
Zebrafish AB9
List of cell lines
| Technology | Biotechnology | null |
1107596 | https://en.wikipedia.org/wiki/Linear%20approximation | Linear approximation | In mathematics, a linear approximation is an approximation of a general function using a linear function (more precisely, an affine function). They are widely used in the method of finite differences to produce first order methods for solving or approximating solutions to equations.
Definition
Given a twice continuously differentiable function of one real variable, Taylor's theorem for the case states that
where is the remainder term. The linear approximation is obtained by dropping the remainder:
This is a good approximation when is close enough to since a curve, when closely observed, will begin to resemble a straight line. Therefore, the expression on the right-hand side is just the equation for the tangent line to the graph of at . For this reason, this process is also called the tangent line approximation. Linear approximations in this case are further improved when the second derivative of a, , is sufficiently small (close to zero) (i.e., at or near an inflection point).
If is concave down in the interval between and , the approximation will be an overestimate (since the derivative is decreasing in that interval). If is concave up, the approximation will be an underestimate.
Linear approximations for vector functions of a vector variable are obtained in the same way, with the derivative at a point replaced by the Jacobian matrix. For example, given a differentiable function with real values, one can approximate for close to by the formula
The right-hand side is the equation of the plane tangent to the graph of at
In the more general case of Banach spaces, one has
where is the Fréchet derivative of at .
Applications
Optics
Gaussian optics is a technique in geometrical optics that describes the behaviour of light rays in optical systems by using the paraxial approximation, in which only rays which make small angles with the optical axis of the system are considered. In this approximation, trigonometric functions can be expressed as linear functions of the angles. Gaussian optics applies to systems in which all the optical surfaces are either flat or are portions of a sphere. In this case, simple explicit formulae can be given for parameters of an imaging system such as focal distance, magnification and brightness, in terms of the geometrical shapes and material properties of the constituent elements.
Period of oscillation
The period of swing of a simple gravity pendulum depends on its length, the local strength of gravity, and to a small extent on the maximum angle that the pendulum swings away from vertical, , called the amplitude. It is independent of the mass of the bob. The true period T of a simple pendulum, the time taken for a complete cycle of an ideal simple gravity pendulum, can be written in several different forms (see pendulum), one example being the infinite series:
where L is the length of the pendulum and g is the local acceleration of gravity.
However, if one takes the linear approximation (i.e. if the amplitude is limited to small swings, ) the period is:
In the linear approximation, the period of swing is approximately the same for different size swings: that is, the period is independent of amplitude. This property, called isochronism, is the reason pendulums are so useful for timekeeping. Successive swings of the pendulum, even if changing in amplitude, take the same amount of time.
Electrical resistivity
The electrical resistivity of most materials changes with temperature. If the temperature T does not vary too much, a linear approximation is typically used:
where is called the temperature coefficient of resistivity, is a fixed reference temperature (usually room temperature), and is the resistivity at temperature . The parameter is an empirical parameter fitted from measurement data. Because the linear approximation is only an approximation, is different for different reference temperatures. For this reason it is usual to specify the temperature that was measured at with a suffix, such as , and the relationship only holds in a range of temperatures around the reference. When the temperature varies over a large temperature range, the linear approximation is inadequate and a more detailed analysis and understanding should be used.
| Mathematics | Basics_2 | null |
1108210 | https://en.wikipedia.org/wiki/Antlion | Antlion | The antlions are a group of about 2,000 species of insect in the neuropteran family Myrmeleontidae. They are known for the predatory habits of their larvae, which mostly dig pits to trap passing ants or other prey. In North America, the larvae are sometimes referred to as doodlebugs because of the marks they leave in the sand. The adult insects are less well known due to their relatively short lifespans in comparison with the larvae. Adults, sometimes known as antlion lacewings, mostly fly at dusk or just after dark and may be mistakenly identified as dragonflies or damselflies.
Antlions have a worldwide distribution. The greatest diversity occurs in the tropics, but a few species are found in cold-temperate locations, one such being the European Euroleon nostras. They most commonly occur in dry and sandy habitats where the larvae can easily excavate their pits, but some larvae hide under debris or ambush their prey among leaf litter.
Antlions are poorly represented in the fossil record. Myrmeleontiformia is generally accepted to be a monophyletic group, and within the Myrmeleontoidea, the antlions' closest living relatives are thought to be the owlflies (Ascalaphidae). A 2019 study finds Myrmeleontidae to be monophyletic, aside from Stilbopteryginae and Palparinae, which form separate clades closer to Ascalaphidae. The predatory actions of the larvae have attracted attention throughout history and antlions have been mentioned in literature since classical times.
Etymology
The exact meaning of the name "antlion" is uncertain. It has been thought to refer to ants forming a large percentage of the prey of the insect, the suffix "lion" merely suggesting "destroyer" or "hunter". In any case, the term seems to go back to classical antiquity. The antlion larva is often called a "doodlebug" in North America because of the odd winding, spiralling trails it leaves in the sand while relocating, which look as if someone has been doodling.
The scientific name of the type genus Myrmeleo – and thus, the family as a whole – is derived from Ancient Greek mýrmex (μύρμηξ) "ant" + léon (λέων) "lion", in a loan translation of the names common across Europe. In most European and Middle Eastern languages, at least the larvae are known under the local term corresponding to "antlion".
Description
Antlions can be fairly small to very large neuropterans, with wingspans ranging from . The African genus Palpares contains some of the largest examples. Acanthaclisis occitanica is the largest European species, with an wingspan, and most North American species approach this size.
The adult has two pairs of long, narrow, multiveined, translucent wings and a long, slender abdomen. Although they somewhat resemble dragonflies or damselflies, they belong to a different infraclass of winged insects. Antlion adults are easily distinguished from damselflies by their prominent, apically clubbed antennae which are about as long as the head and thorax combined. Also, the pattern of wing venation differs, and compared to damselflies, the adults are very feeble fliers and are normally found fluttering about at night in search of a mate. Adult antlions are typically nocturnal, and rarely seen by day.
Males of most species have a unique structure, a bristle-bearing knob known as a "pilula axillaris", at the base of the rear wing. The abdomen in males is usually longer than in females and often has an extra lobe. The tip of the abdomen of females shows greater variation than that of males, depending perhaps on oviposition sites, and usually bears tufts of bristles for digging and a finger-like extension.
The antlion larva has a robust fusiform body, a very plump abdomen, and a thorax bearing three pairs of walking legs. The prothorax forms a slender mobile "neck" for the large, square, flattened head, which bears an enormous pair of sickle-like jaws with several sharp, hollow projections. The jaws are formed by the maxillae and mandibles; the mandibles each contain a deep groove over which the maxilla fits neatly, forming an enclosed canal for injecting venom to immobilise the victim, and enzymes to digest its soft parts. The larva is clad in forward-pointing bristles which help it to anchor itself and exert greater traction, enabling it to subdue prey considerably larger than itself.
Antlion larvae are unusual among insects in lacking an anus. All the metabolic waste generated during the larval stage is stored; some is used to spin the silk for the cocoon and the rest is eventually voided as meconium at the end of its pupal stage.
Distribution
There are about 2,000 species of antlion found in most parts of the world, with the greatest diversity being in warmer areas. The best known species are those in which the larvae dig pits to trap their prey, but not all species do this. Antlions live in a range of usually dry habitats including open woodland floors, scrub-clad dunes, hedge bases, river banks, road verges, under raised buildings and in vacant lots.
Life-cycle
Apart from pit-trap-forming taxa, the biology of members of the family Myrmeleontidae, to which the antlions belong, has been little studied. The life-cycle begins with oviposition (egg-laying) in a suitable location. The female antlion repeatedly taps the prospective laying site with the tip of her abdomen and then inserts her ovipositor into the substrate and lays an egg.
Depending on the species and where it lives, the larva either conceals itself under leaves, debris or pieces of wood, hides in a crack or digs a funnel-shaped pit in loose material. As ambush predators, catching prey is risky because food arrives unpredictably and, for those species that make traps, maintaining one is costly. The larvae therefore have low metabolic rates and can survive for long periods without food. They can take several years to complete their life-cycle; they mature faster with plentiful food, but can survive for many months without feeding. In cooler climates they dig their way deeper and remain inactive during the winter.
When the larva attains its maximum size, it pupates and undergoes metamorphosis. It makes a globular cocoon of sand or other local substrate stuck together with fine silk spun from a slender spinneret at the rear end of the body. The cocoon may be buried several centimetres deep in sand. After completing its transformation into an adult insect over the course of about one month, it emerges from the case, leaving the pupal integument behind, and works its way to the surface. After about twenty minutes, the adult's wings are fully opened and it flies off in search of a mate. The adult is considerably larger than the larva as antlions exhibit the greatest disparity in size between larva and adult of any type of holometabolous insect. This is by virtue of the fact that the exoskeleton of the adult is extremely thin and flimsy, with an exceptionally low density. The adult typically lives for about 25 days, but some insects survive for as long as 45 days.
Ecology
Antlion larvae eat small arthropods – mainly ants – while the adults of some species eat pollen and nectar, and others are predators of small arthropods. In certain species of Myrmeleontidae, such as Dendroleon pantherinus, the larva, although resembling that of Myrmeleon structurally, makes no pitfall trap, but hides in detritus in a hole in a tree and seizes passing prey. In Japan, Gatzara jezoensis larvae lurk on the surface of rocks for several years while awaiting prey; during this time they often become coated with lichen, and have been recorded at densities of up to 344 per square metre.
The larva is a voracious predator. Within a few minutes of seizing its prey with its jaws and injecting it with venom and enzymes, it begins to suck out the digestion products. The larva is extremely sensitive to ground vibrations, the low-frequency sounds made by an insect crawling across the ground; the larva locates the source of the vibrations by the differences in timing of the arrival of waves detected by receptors, tufts of hairs on the sides of the two hindmost thoracic segments.
Pit-building species
Funnel-shaped pits are built by members of just 3 antlion tribes: Myrmeleontini, Myrmecaelurini, and Nesoleontini. In these trap-building species, an average-sized larva digs a pit about 2 in (5 cm) deep and 3 in (7.5 cm) wide at the edge. This behavior has also been observed in the Vermileonidae (Diptera), whose larvae dig the same sort of pit to feed on ants. Having marked out the chosen site with a circular groove, the antlion larva begins to crawl backward, using its abdomen as a plough to dig up the soil. Using one of its front legs, it places heaps of loosened particles upon its head which it then flicks clear of the scene of operations. Continuing in this way, it gradually works its way from the circumference toward the center. As it slowly moves round and round, the pit gradually gets deeper and deeper, until the slope angle reaches the critical angle of repose (that is, the steepest angle the sand can maintain, where it is on the verge of collapse from slight disturbance), and the pit is solely lined by fine grains. By digging in a spiral when constructing its pit, the antlion minimises the time needed to complete the pit.
When the pit is completed, the larva settles down at the bottom, buried in the soil with only the jaws projecting above the surface, often in a wide-opened position on either side of the very tip of the cone. The steep-sloped trap that guides prey into the larva's mouth while avoiding crater avalanches is one of the simplest and most efficient traps in the animal kingdom. The fine grain lining ensures that the avalanches which carry prey are as large as possible.
Since the sides of the pit consist of loose sand at its angle of repose, they afford an insecure foothold to any small insects that inadvertently venture over the edge, such as ants. Slipping to the bottom, the prey is immediately seized by the lurking antlion; if it attempts to scramble up the treacherous walls of the pit, it is speedily checked in its efforts and brought down by showers of loose sand which are thrown at it from below by the larva. By throwing up loose sand from the bottom of the pit, the larva also undermines the sides of the pit, causing them to collapse and bring the prey with them. Thus, it does not matter whether the larva actually strikes the prey with the sand showers.
Antlion larvae are capable of capturing and killing a variety of insects and other arthropods, and can even subdue small spiders. The projections in the jaws of the larva are hollow and through this, the larva sucks the fluids out of its victim. After the contents are consumed, the dry carcass is flicked out of the pit. The larva readies the pit once again by throwing out collapsed material from the center, steepening the pit walls to the angle of repose.
Antlion larvae require loose soil, not necessarily, but often, sand. Antlions can also handle larger granular material which is filtered out of the soil during pit construction. The larvae prefer dry places protected from the rain. When it first hatches, the tiny larva specialises in very small insects, but as it grows larger, it constructs larger pits, and thus catches larger prey, sometimes much larger than itself.
Other arthropods may make use of the antlion larva's ability to trap prey. The larva of the Australian horsefly (Scaptia muscula) lives in antlion (for example Myrmeleon pictifrons) pit traps and feeds on the prey caught, and the female chalcid wasp (Lasiochalcidia igiliensis) purposefully allows itself to be trapped so that it can parasitise the antlion larva by ovipositing between its head and thorax.
Recent research has found that antlion larvae often "play dead" for a variable amount of time (from a few minutes up to an hour) when disturbed to hide from predators. The method is effective; it increased survival rates in patches that use it by 20%. Furthermore, they appear to have maximized its usefulness—further increasing the duration is not likely to convey substantial survival benefits to the larvae.
Evolution
The closest living relatives of antlions within the Myrmeleontoidea are the owlflies (Ascalaphidae); the Nymphidae are more distantly related. The extinct Araripeneuridae and Babinskaiidae are considered likely to be stem groups in the Myrmeleontiformia clade. The phylogeny of the Neuroptera has been explored using mitochondrial DNA sequences, and while issues remain for the group as a whole (the "Hemerobiiformia" being paraphyletic), the Myrmeleontiformia is generally agreed to be monophyletic, giving the following cladogram:
The subfamilies are shown below; a few genera, mostly fossil, are of uncertain or basal position. The fossil record of antlions is very small by neuropteran standards. However, some Mesozoic fossils attest to the antlions' origin more than 150 million years ago. These were at one time separated as the Palaeoleontidae, but are now usually recognized as early antlions.
Taxonomy
The supra-generic classification within the Myrmeleontidae is disputed by different authors. Several different schematics have been proposed.
Stange (2004) classification
Stange recognized three subfamilies, consisting of Stilbopteryginae, Palparinae, and Myrmeleontinae. Brachynemurinae, Dendroleontinae, and other subfamilies recognized by prior authors were placed in the Myrmeleontinae.
Stilbopteryginae Newman, 1853
Stilbopterygini monotypic: Stilbopteryx
Palparinae Banks, 1911
Dimarini Navas, 1914
Palparidiini Markl, 1954
Palparini Banks, 1911
Pseudimarini Markl 1954
Myrmeleontinae Latreille, 1802
Acanthaclisini Navas, 1912
Brachynemurini Banks, 1927
Dendroleontini Banks, 1899
Gnopholeontini Stange, 1994
Lemolemini Stange, 1994
Maulini Markl 1954
Myrmecaelurini Esben-Petersen, 1918
Myrmeleontini Latreille, 1802
Nesoleontini Markl, 1954
Nemoleontini Banks, 1911
Michel et al. (2017) classification
Michel et.al recognized just four subfamilies, consisting of Acanthaclisinae, Myrmeleontinae, Palparinae, and Stilbopteryginae.
Stilbopteryginae Newman, 1853
Stilbopterygini Newman, 1853
Acanthaclisinae Navas, 1912
Acanthaclisini Navas, 1912
Palparinae Banks, 1911
Dimarini Navas, 1914
Palparidiini Markl, 1954
Palparini Banks, 1911
Pseudimarini Markl 1954
Myrmeleontinae Latreille, 1802
Brachynemurini Banks, 1927
Dendroleontini Banks, 1899
Gepini Markl, 1954
Gnopholeontini Stange, 1994
Lemolemini Stange, 1994
Maulini Markl 1954
Myrmecaelurini Esben-Petersen, 1918
Myrmeleontini Latreille, 1802
Nesoleontini Markl, 1954
Machado et al. (2018) classification
A subsequent revision by Machado et.al recognized a different four subfamilies, notable in the inclusion of the family Ascalaphidae as a subfamily instead of as a sister taxon to Myrmeleontidae: Ascalaphinae, Myrmeleontinae, Dendroleontinae, and Nemoleontinae.
Ascalaphinae Rambur, 1842
Ascalaphini Lefèbvre, 1842
Dimarini Navas, 1914
Haplogleniini Newman, 1853
Palparini Banks, 1911
Stilbopterygini Newman, 1853
Ululodini Van der Weele 1909
Myrmeleontinae Latreille, 1802
Acanthaclisini Navas, 1912
Brachynemurini Banks, 1927
Myrmecaelurini Esben-Petersen, 1918
Myrmeleontini Latreille, 1802
Nesoleontini Markl, 1954
Dendroleontinae Banks, 1899
Acanthoplectrini Markl, 1954
Dendroleontini Banks, 1899
Glenurini Banks, 1927
Megistopini Navas, 1912
Nemoleontinae Banks, 1911
Nemoleontini Banks, 1911
Protoplectrini Tillyard, 1916
Jones (2019) classification
Jones dissented with Machado et al. and presented an alternative classification, restoring the Ascalaphidae to its traditional family-rank placement, and elevating Palparinae and Stilbopteryginae to family level, leaving only a single subfamily in Myrmeleontidae sensu stricto. However, this classification has been rejected by Hévin et al. (2023), finding the taxonomic decisions to be "not substantiated, with only the family name Palparidae being mentioned in isolation."
Myrmeleontinae Latreille, 1802
Acanthaclisini Navas, 1912
Brachynemurini Banks, 1927
Dendroleontini Banks, 1899
Myrmecaelurini Esben-Petersen, 1918
Myrmeleontini Latreille, 1802
Nemoleontini Banks, 1911
Nesoleontini Markl, 1954
In culture and folklore
In popular folklore in the southern United States, people recite a poem or chant to make the antlion come out of its hole. Similar practices have been recorded from Africa, the Caribbean, China and Australia.
The Myrmecoleon was a mythical ant–lion hybrid written about in the 2nd century AD Physiologus, where animal descriptions were paired with Christian morals. The ant-lion as described was said to starve to death because of its dual nature – the lion nature of the father could only eat meat, but the ant half from the mother could only eat grain chaff, thus the offspring could not eat either and would starve. It was paired with the Biblical verse Matthew 5:37. The fictional ant-lion of Physiologus is probably derived from a misreading of Job 4:11.
The French naturalist Jean-Henri Fabre wrote that "The Ant-lion makes a slanting funnel in the sand. Its victim, the Ant, slides down the slant and is then stoned, from the bottom of the funnel, by the hunter, who turns his neck into a catapult."
In Japan, both the insect and its pit-traps are popularly known as . This term has since become a stock colloquialism for any "inescapable" trap, whether literal or metaphorical (i.e. an unpleasant social obligation).
Antlions appear as antagonists in the 1991 life simulation video game, SimAnt, and (in a giant form) in the Final Fantasy series, Grounded, Terraria, Don't Starve Together, Monster Rancher 2, Mother 3 and in the Half-Life 2 video game series as an unrelated alien insect species sharing sand burowing traits with the real antlion larvae. The Trapinch, Vibrava, and Flygon Pokémon evolution line is based on an antlion. The fictional sarlacc from the Star Wars franchise is often compared to the real-life antlion. It also appears as a predator in the film Enemy Mine.
In the third book of Tove Jansson's Moomins series, Finn Family Moomintroll, a rather large and fanciful antlion appears in the second chapter, depicted as a sand-dwelling predator with the literal head of a lion.
| Biology and health sciences | Insects: General | Animals |
1108628 | https://en.wikipedia.org/wiki/Homo%20floresiensis | Homo floresiensis | Homo floresiensis , also known as "Flores Man" or "Hobbit" (after the fictional species), is an extinct species of small archaic humans that inhabited the island of Flores, Indonesia, until the arrival of modern humans about 50,000 years ago.
The remains of an individual who would have stood about in height were discovered in 2003 at Liang Bua cave. As of 2015, partial skeletons of fifteen individuals have been recovered, including one complete skull, referred to as "LB1".
Homo floresiensis is thought to have arrived on Flores around 1.27–1 million years ago. There is debate as to whether H. floresiensis represents a descendant of Javanese Homo erectus that reduced its body size as a result of insular dwarfism, or whether it represents an otherwise undetected migration of small, Australopithecus or Homo habilis-grade archaic humans outside of Africa.
This hominin was at first considered remarkable for its survival until relatively recent times, initially thought to be only 12,000 years ago. However, more extensive stratigraphic and chronological work has pushed the dating of the most recent evidence of its existence back to 50,000 years ago. The Homo floresiensis skeletal material at Liang Bua is now dated from 60,000 to 100,000 years ago; stone tools recovered alongside the skeletal remains were from archaeological horizons ranging from 50,000 to 190,000 years ago. Other earlier remains from Mata Menge date to around 700,000 years ago.
Specimens
Discovery
The first specimens were discovered on the Indonesian island of Flores on 2 September 2003 by a joint Australian-Indonesian team of archaeologists looking for evidence of the original human migration of modern humans from Asia to Australia. They instead recovered a nearly complete, small-statured skeleton, LB1, in the Liang Bua cave, and subsequent excavations in 2003 and 2004 recovered seven additional skeletons, initially dated from 38,000 to 13,000 years ago.
In 2004, a separate species Homo floresiensis was named and described by Peter Brown et al., with LB1 as the holotype. A tooth, LB2, was referred to the species. LB1 is a fairly complete skeleton, including a nearly complete skull, which belonged to a 30-year-old woman, and has been nicknamed "Little Lady of Flores" or "Flo". An arm bone provisionally assigned to H. floresiensis, specimen LB3, is about 74,000 years old. The specimens are not fossilized and have been described as having "the consistency of wet blotting paper". Once exposed, the bones had to be left to dry before they could be dug up. The discoverers proposed that a variety of features, both primitive and derived, identify these individuals as belonging to a new species. Based on previous date estimates, the discoverers also proposed that H. floresiensis lived contemporaneously with modern humans on Flores. Before publication, the discoverers were considering placing LB1 into her own genus, Sundanthropus floresianus (), but reviewers of the article recommended that, despite her size, she should be placed in the genus Homo.
In 2009, additional finds were reported, increasing the minimum number of individuals represented by bones to fourteen. In 2015, teeth were referred to a fifteenth individual, LB15.
Stone implements of a size considered appropriate to these small humans are also widely present in the cave. The implements are at horizons initially dated to 95,000 to 13,000 years ago. Modern humans reached the region by around 50,000 years ago, by which time H. floresiensis is thought to have gone extinct. Comparisons of the stone artifacts with those made by modern humans in East Timor indicate many technological similarities.
Scandal over specimen damage
The fossils are property of the Indonesian state. In early December 2004, Indonesian paleoanthropologist Teuku Jacob, formerly chief paleontologist of the Indonesian Gadjah Mada University, removed most of the remains from their repository, Jakarta's National Research Centre of Archaeology, with the permission of one of the institute's directors, Raden Panji Soejono, and kept them for three months. Professor Jacob did not believe the specimens represented a different species, contending that the LB1 find was from a 25–30 year-old omnivorous subspecies of H. sapiens, probably a pygmy, and that the small skull was due to microcephaly, which produces a small brain and skull. Professor Richard Roberts of the University of Wollongong in Australia and other anthropologists expressed the fear that important scientific evidence would be sequestered by a small group of scientists who neither allowed access by other scientists nor published their own research. Jacob returned the remains on 23 February 2005 with portions severely damaged and missing two leg bones.
Press reports thus described the condition of the returned remains: "[including] long, deep cuts marking the lower edge of the Hobbit's jaw on both sides, said to be caused by a knife used to cut away the rubber mould ... the chin of a second Hobbit jaw was snapped off and glued back together. Whoever was responsible misaligned the pieces and put them at an incorrect angle ... The pelvis was smashed, destroying details that reveal body shape, gait and evolutionary history.", causing the discovery team leader Morwood to remark, "It's sickening; Jacob was greedy and acted totally irresponsibly."
Jacob, however, denied any wrongdoing. He stated that the damages occurred during transport from Yogyakarta back to Jakarta despite the claimed physical evidence that the jawbone had been broken while making a mould of the bones.
In 2005, Indonesian officials forbade access to the cave. Some news media, such as the BBC, expressed the opinion that the restriction was to protect Jacob, who was considered "Indonesia's king of palaeoanthropology", from being proved wrong. Scientists were allowed to return to the cave in 2007, shortly after Jacob's death.
Classification and evolution
Phylogeny and evolution
Because of the deep neighbouring Lombok Strait, Flores remained an isolated island during episodes of low sea level. Therefore, the ancestors of H. floresiensis could only have reached the island by oceanic dispersal, most likely by rafting. The oldest stone tools on Flores are around 1 million years old. Stone artifacts are absent from sites over 1.27 million years old, suggesting that the ancestors of H. floresiensis arrived after this time.
In 2016, fossil teeth and a partial jaw from hominins assumed to be ancestral to H. floresiensis were discovered at Mata Menge, about from Liang Bua. They date to about 700,000 years ago. Other remains including a humerus were later described from Mata Menge, with the remains subsequently being directly assigned to H. floresiensis. These remains are about the same size or somewhat smaller than the remains from Liang Bua, suggesting the size of the species remained stable for hundreds of thousands of years up until its extinction.
Two hypotheses have been proposed as to the origin of H. floresiensis. The first proposes that H. floresiensis descended from an early migration of very primitive small Australopithecus/Homo habilis-grade archaic humans outside of Africa prior to 1.75 million years ago. This is based on various aspects of H. floresiensis sketetal anatomy, such as its feet bones being considered as more similar to those of very archaic humans such as Australopithecus and Homo habilis than to Homo erectus. This position has been supported by several cladistical analyses. Other authors have argued that H. floresiensis instead likely represents the descendants of a population of Javanese Homo erectus that became isolated on Flores, with the small body size being the result of insular dwarfism, a well known evolutionary trend found among various island animals. These authors alternatively suggest that H. floresiensis has several cranial and dental similarities to H. erectus, particularly to early Javanese Homo erectus. These authors also dispute some of the similarities to Australopithecus and Homo habilis-grade archaic humans, and suggest that others may have been the result of evolutionary reversals/convergence.
It has been noted that there is no evidence archaic humans in the adjacent (and likely source) region of Java earlier than 1.3-1.5 or 1.8 million years ago, with the earliest human presence on Java being represented by Homo erectus, with there also being no evidence of Australopithecus or Homo habilis-grade archaic humans anywhere outside of Africa, which supporters of the Homo erectus-origin hypothesis suggest makes the descent of H. floresiensis from these more primitive hominins unlikely.
DNA extraction attempt
In 2006, two teams attempted to extract DNA from a tooth discovered in 2003, but both teams were unsuccessful. It has been suggested that this happened because the dentine was targeted; new research suggests that the cementum has higher concentrations of DNA. Moreover, the heat generated by the high speed of the drill bit may have denatured the DNA.
Congenital disorder claims
The small brain size of H. floresiensis at 417 cc prompted hypotheses that the specimens were simply H. sapiens with a birth defect, rather than the result of neurological reorganisation. These claims have subsequently been widely rejected.
Microcephaly
Prior to Jacob's removal of the fossils, American neuroanthropologist Dean Falk and her colleagues performed a CT scan of the LB1 skull and a virtual endocast, and concluded that the brainpan was neither that of a pygmy nor an individual with a malformed skull and brain. In response, American neurologist Jochen Weber and colleagues compared the computer model skull with microcephalic human skulls, and found that the skull size of LB1 falls in the middle of the size range of the human samples, and is not inconsistent with microcephaly. A 2006 study stated that LB1 probably descended from a pygmy population of modern humans, but herself shows signs of microcephaly, and other specimens from the cave show small stature but not microcephaly.
In 2005, the original discoverers of H. floresiensis, after unearthing more specimens, countered that the skeptics had mistakenly attributed the height of H. floresiensis to microcephaly. Falk stated that Martin's assertions were unsubstantiated. In 2006, Australian palaeoanthropologist Debbie Argue and colleagues also concluded that the finds are indeed a new species. In 2007, Falk found that H. floresiensis brains were similar in shape to modern humans, and the frontal and temporal lobes were well-developed, which would not have been the case were they microcephalic.
In 2008, Greek palaeontologist George Lyras and colleagues said that LB1 falls outside the range of variation for human microcephalic skulls. However, a 2013 comparison of the LB1 endocast to a set of 100 normocephalic and 17 microcephalic endocasts showed that there is a wide variation in microcephalic brain shape ratios and that in these ratios the group as such is not clearly distinct from normocephalics. The LB1 brain shape nevertheless aligns slightly better with the microcephalic sample, with the shape at the extreme edge of the normocephalic group. A 2016 pathological analysis of LB1's skull revealed no pathologies nor evidence of microcephaly, and concluded that LB1 is a separate species.
Laron syndrome
A 2007 study postulated that the skeletons were those of humans who suffered from Laron syndrome, which was first reported in 1966, and is most common in inbreeding populations, which may have been the scenario on the small island. It causes a short stature and small skull, and many conditions seen in Laron syndrome patients are also exhibited in H. floresiensis. The estimated height of LB1 is at the lower end of the average for afflicted human women, but the endocranial volume is much smaller than anything exhibited in Laron syndrome patients. DNA analysis would be required to support this theory.
Congenital iodine deficiency syndrome
In 2008 Australian researcher Peter Obendorf — who studies congenital iodine deficiency syndrome — and colleagues suggested that LB1 and LB6 suffered from myxoedematous (ME) congenital iodine deficiency syndrome resulting from congenital hypothyroidism (underactive thyroid), and that they were part of an affected population of H. sapiens on the island. Congenital iodine deficiency syndrome, caused by iodine deficiency, is expressed by small bodies and reduced brain size (but ME causes less motor and mental disablement than other forms of congenital iodine deficiency syndrome), and is a form of dwarfism still found in the local Indonesian population. They said that various features of H. floresiensis are diagnostic characteristics, such as enlarged pituitary fossa, unusually straight and untwisted humeral heads, relatively thick limbs, double rooted premolar, and primitive wrist morphology.
However, Falk's scans of LB1's pituitary fossa show that it is not larger than usual. Also, in 2009, anthropologists Colin Groves and Catharine FitzGerald compared the Flores bones with those of ten people who had had cretinism, and found no overlap. Obendorf and colleagues rejected Groves and FitzGerald's argument the following year. A 2012 study similar to Groves and FitzGeralds' also found no evidence of congenital iodine deficiency syndrome.
Down syndrome
From 2006, physical anthropologist Maciej Henneberg and colleagues have claimed that LB1 suffered from Down syndrome, and that the remains of other individuals at the Flores site were merely normal modern humans. However, there are a number of characteristics shared by both LB1 and LB6 as well as other known early humans and absent in H. sapiens, such as the lack of a chin. In 2016, a comparative study concluded that LB1 did not exhibit a sufficient number of Down syndrome characteristics to support a diagnosis. The noted physical anthropologist Chris Stringer in 2011 wrote of Homo floresiensis deniers generally, "I think they have damaged their own, and palaeoanthropologist's, reputation."
Anatomy
The most important and obvious identifying features of Homo floresiensis are its small body and small cranial capacity. Brown and Morwood also identified a number of additional, less obvious features that might distinguish LB1 from modern H. sapiens, including the form of the teeth, the absence of a chin, and a lesser torsion in the lower end of the humerus (upper arm bone). Each of these putative distinguishing features has been heavily scrutinized by the scientific community, with different research groups reaching differing conclusions as to whether these features support the original designation of a new species, or whether they identify LB1 as a severely pathological H. sapiens.
A 2015 study of the dental morphology of forty teeth of H. floresiensis compared to 450 teeth of living and extinct human species, states that they had "primitive canine-premolar and advanced molar morphologies," which is unique among hominins.
The discovery of additional partial skeletons has verified the existence of some features found in LB1, such as the lack of a chin, but Jacob and other research teams argue that these features do not distinguish LB1 from local modern humans. Lyras et al. have asserted, based on 3D-morphometrics, that the skull of LB1 differs significantly from all H. sapiens skulls, including those of small-bodied individuals and microcephalics, and is more similar to the skull of Homo erectus. Ian Tattersall argues that the species is wrongly classified as Homo floresiensis as it is far too archaic to assign to the genus Homo.
Size
LB1's height is estimated to have been . The height of a second skeleton, LB8, has been estimated at based on tibial length. These estimates are outside the range of normal modern human height and considerably shorter than the average adult height of even the smallest modern humans, such as the Mbenga and Mbuti at , Twa, Semang at for adult women of the Malay Peninsula, or the Andamanese at also for adult women. LB1's body mass is estimated to have been . LB1 and LB8 are also somewhat smaller than the australopithecines, such as Lucy, from three million years ago, not previously thought to have expanded beyond Africa. Thus, LB1 and LB8 may be the shortest and smallest members of the extended human group discovered thus far.
Their short stature was likely due to insular dwarfism, where size decreases as a response to fewer resources in an island ecosystem. In 2006, Indonesian palaeoanthropologist Teuku Jacob and colleagues said that LB1 has a similar stature to the Rampasasa pygmies who inhabit the island, and that size can vary substantially in pygmy populations. A 2018 study refuted the possibility of Rampasasa pygmies descending from H. floresiensis, concluding that "multiple independent instances of hominin insular dwarfism occurred on Flores".
Aside from smaller body size, the specimens seem to otherwise resemble H. erectus, a species known to have been living in Southeast Asia at times coincident with earlier finds purported to be of H. floresiensis.
Brain
In addition to a small body size, H. floresiensis had a remarkably small brain size. LB1's brain is estimated to have had a volume of , placing it at the range of chimpanzees or the extinct australopithecines. LB1's brain size is less than half that of its presumed immediate ancestor, H. erectus (). The brain-to-body mass ratio of LB1 lies between that of H. erectus and the great apes.<ref name="Falk_et_al_2006">{{cite journal| author = Falk, D.|display-authors=etal |date= 19 May 2006 |title= Response to Comment on 'The Brain of LB1, Homo floresiensis''' |journal= Science |volume= 312 |issue= 5776 |pages= 999c |bibcode = 2006Sci...312.....F |doi= 10.1126/science.1124972 |doi-access= free }}</ref> Such a reduction is likely due to insular dwarfism, and a 2009 study found that the reduction in brain size of extinct pygmy hippopotamuses in Madagascar compared with their living relatives is proportionally greater than the reduction in body size, and similar to the reduction in brain size of H. floresiensis compared with H. erectus.
Smaller size does not appear to have affected mental faculties, as Brodmann area 10 on the prefrontal cortex, which is associated with cognition, is about the same size as that of modern humans. H. floresiensis is also associated with evidence for advanced behaviours, such as the use of fire, butchering, and stone tool manufacturing.
Limbs
The angle of humeral torsion was lesser than in modern humans. The humeral head of modern humans is twisted between 145 and 165 degrees to the plane of the elbow joint, whereas it is 120 degrees in H. floresiensis. This may have provided an advantage when arm-swinging, and, in tandem with the unusual morphology of the shoulder girdle and short clavicle, would have displaced the shoulders slightly forward into an almost shrugging position. The shrugging position would have compensated for the lower range of motion in the arm, allowing for similar maneuverability in the elbows as in modern humans. The wrist bones are similar to those of apes and Australopithecus. They were significantly smaller and more flexible than the carpals of modern humans, lacking contemporary features which evolved at least 800,000 years ago.
The leg bones were more robust than those of modern humans. The feet were unusually flat and large in relation with the rest of the body. As a result, when walking, they would have had to bend the knees further back than modern humans do. This caused a high-stepping gait and slow walking speed. The toes were thin, long, and halluces were almost indistinguishable from the other metatarsals.
Culture
The cave yielded over ten thousand stone artefacts, mainly lithic flakes, surprising considering H. floresiensiss small brain. This has led some researchers to theorize that H. floresiensis inherited their tool-making skills from H. erectus. Points, perforators, blades, and microblades were associated with remains of the extinct elephant-relative Stegodon. It has therefore been proposed that H. floresiensis hunted juvenile Stegodon. Similar artefacts are found at the Soa Basin south, associated with Stegodon and Komodo dragon remains, and are attributed to a likely ancestral population of H. erectus. Other authors have doubted the extent of hunting of Stegodon by H. floresiensis, noting the rarity of cut marks on remains of Stegodon found at Liang Bua, suggesting that they would have faced intense competition for carcasses with other predators, like the Komodo dragon, the giant stork Leptoptilos robustus, and vultures, and that it was possible that their main prey was instead the giant rats like Papagomys endemic to the island, which are found abundantly at Liang Bua. While it was initially suggested that H. floresiensis was capable of using fire, the supporting evidence for this claim was later found to be unreliable.
Extinction
The youngest H. floresiensis bone remains in the cave date to 60,000 years ago, and the youngest stone tools to 50,000 years ago. The previous estimate of 12,000 BP was due to an undetected unconformity in the cave stratigraphy. The timing of their disappearance from the cave stratigraphy is close to the time that modern humans reached the area, which may suggest the effects of modern humans directly on H. floresiensis or more broadly on the ecosystems of Flores caused or contributed to their extinction. DNA analysis of pygmy modern humans from Flores has found no evidence of any DNA from H. floresiensis.
Paleoecology
During the late Early Pleistocene-Late Pleistocene before the arrival of Homo sapiens, Flores exhibited a depauperate ecosystem with relatively few terrestrial vertebrate species, including the extinct dwarf proboscidean (elephant relative) Stegodon florensis; and a variety of rats (Murinae) including small-sized forms like Rattus hainaldi, the Polynesian rat, Paulamys, and Komodomys, the medium-sized Hooijeromys, and giant Papagomys and extinct Spelaeomys, the latter two genera being about the size of rabbits, with body masses of . Also present were the Komodo dragon and another smaller monitor lizard (Varanus hooijeri), with birds including a giant stork (Leptoptilos robustus) and a vulture (Trigonoceps).
"Hobbit" nicknameHomo floresiensis was swiftly nicknamed "the hobbit" by the discoverers, after the fictional race popularized in J. R. R. Tolkien's book The Hobbit, and some of the discoverers suggested naming the species H. hobbitus.
In October 2012, a New Zealand scientist due to give a public lecture on Homo floresiensis was told by the Tolkien Estate that he was not allowed to use the word "hobbit" in promoting the lecture.
In 2012, the American film studio The Asylum, which produces low-budget "mockbuster" films, planned to release a movie entitled Age of the Hobbits depicting a "peace-loving" community of H.floresiensis "enslaved by the Java Men, a race of flesh-eating dragon-riders." The film was intended to piggyback on the success of Peter Jackson's film The Hobbit: An Unexpected Journey. The film was blocked from release due to a legal dispute about using the word "hobbit." The Asylum argued that the film did not violate the Tolkien copyright because the film was about H.floresiensis, "uniformly referred to as 'Hobbits' in the scientific community." The film was later retitled Clash of the Empires''.
| Biology and health sciences | Homo | Biology |
1109384 | https://en.wikipedia.org/wiki/Soxhlet%20extractor | Soxhlet extractor | A Soxhlet extractor is a piece of laboratory apparatus invented in 1879 by Franz von Soxhlet. It was originally designed for the extraction of a lipid from a solid material. Typically, Soxhlet extraction is used when the desired compound has a limited solubility in a solvent, and the impurity is insoluble in that solvent. It allows for unmonitored and unmanaged operation while efficiently recycling a small amount of solvent to dissolve a larger amount of material.
Description
A Soxhlet extractor has three main sections: a percolator (boiler and reflux) which circulates the solvent, a thimble (usually made of thick filter paper) which retains the solid to be extracted, and a siphon mechanism, which periodically empties the condensed solvent from the thimble back into the percolator.
Assembly
The source material containing the compound to be extracted is placed inside the thimble.
The thimble is loaded into the main chamber of the Soxhlet extractor.
The extraction solvent to be used is placed in a distillation flask.
The flask is placed on the heating element.
The Soxhlet extractor is placed atop the flask.
A reflux condenser is placed atop the extractor.
Operation
The solvent is heated to reflux. The solvent vapour travels up a distillation arm, and floods into the chamber housing the thimble of solid. The condenser ensures that any solvent vapour cools, and drips back down into the chamber housing the solid material. The chamber containing the solid material slowly fills with warm solvent. Some of the desired compound dissolves in the warm solvent. When the Soxhlet chamber is almost full, the chamber is emptied by the siphon. The solvent is returned to the distillation flask. The thimble ensures that the rapid motion of the solvent does not transport any solid material to the still pot. This cycle may be allowed to repeat many times, over hours or days.
During each cycle, a portion of the non-volatile compound dissolves in the solvent. After many cycles the desired compound is concentrated in the distillation flask. The advantage of this system is that instead of many portions of warm solvent being passed through the sample, just one batch of solvent is recycled.
After extraction the solvent is removed, typically by means of a rotary evaporator, yielding the extracted compound. The non-soluble portion of the extracted solid remains in the thimble, and is usually discarded.
Like Soxhlet extractor, the Kumagawa extractor has a specific design where the thimble holder/chamber is directly suspended inside the solvent flask (having a vertical large opening) above the boiling solvent. The thimble is surrounded by hot solvent vapour and maintained at a higher temperature compared to the Soxhlet extractor, thus allowing better extraction for compounds with higher melting points such as bitumen. The removable holder/chamber is fitted with a small siphon side arm and, in the same way as for Soxhlet, a vertical condenser ensures that the solvent drips back down into the chamber which is automatically emptied at every cycle.
History
William B. Jensen notes that the earliest example of a continuous extractor is archaeological evidence for a Mesopotamian hot-water extractor for organic matter dating from approximately 3500 BC. The same mechanism is present in the Pythagorean cup. Before Soxhlet, the French chemist Anselme Payen also pioneered with continuous extraction in the 1830s.
A Soxhlet apparatus has been proposed as an effective technique for washing mass standards.
Gallery
| Physical sciences | Other separations | Chemistry |
4388926 | https://en.wikipedia.org/wiki/Snipe%20eel | Snipe eel | Snipe eels are a family, Nemichthyidae, of eels that consists of nine species in three genera. They are pelagic fishes, found in every ocean, mostly at depths of but sometimes as deep as . Depending on the species, adults may reach in length, yet they weigh only . They are distinguished by their very slender jaws that separate toward the tips as the upper jaw curves upward. The jaws appear similar to the beak of the bird called the snipe. Snipe eels are oviparous, and the juveniles, called Leptocephali (meaning small head), do not resemble the adults but have oval, leaf-shaped and transparent bodies. Different species of snipe eel have different shapes, sizes and colors. The similarly named bobtail snipe eel is actually in a different family and represented by two species, the black Cyema atrum and the bright red Neocyema erythrosoma.
Genera and species
There are nine species in three genera:
Characteristics
Snipe eels have the long, slender body typical of eels. They are distinguished by their very slender jaws that separate toward the tips as the upper jaw curves upward and which appear similar to the beak of the bird called the snipe. The vent may be near the throat (as in N. scolopaceus) or further back on the body (as in A. gilli and A. infans). Depending on the species, adults may reach 30–60 inches long, and weigh only a few ounces to a pound. The dorsal fins show wide variation, too. For example, A. gilli has fewer than 280 rays in its dorsal fin, while A. infans has more than 320.
The only food items actually found in the stomachs of snipe eels have been shrimp-like crustaceans, though ichthyologists believe they should be capable of catching and eating small fish and cephalopods also. Since predatory fish often feed on eel larvae, it is presumed by scientists and ichthyologists that they feed on snipe eel larvae as well. It is also thought that larger fish eat adult snipe eels though there is not much direct evidence for this.
Distribution
Snipe eels are found in every ocean and generally occupy depths of 300–600 m, though specimens have been caught nearer the surface at night, and storms occasionally result in individuals being stranded on the shore. Larval snipe eels occupy more shallow regions of 60–70 m before descending to metamorphose into adult form.
Reproduction
Based on studies of several species of snipe eel in the Sargasso Sea and off the coast of California, it seems that snipe eels spawn mainly in the spring, but also into early summer. The juvenile forms of eels are generally called leptocephali (meaning "small head") and do not look like their adult forms. They are flat and transparent, taking forms that may be elongated or closer to the shape of a leaf. They remain in the larval form and near the surface (upper 200 m) for several months before descending to metamorphose. This is similar to other eel species but much longer than most other fish.
Similar fish
The bobtail snipe eels are two species of deep-sea fishes in the family Cyematidae, one only in each of two genera. They are small elongate fishes, growing up to long. They are bathypelagic (deep-water ocean-dwellers) and have been found down to . They are found in all oceans. Cyema atrum is black, while Neocyema erythrosoma is bright red.
Sawtooth eels are a family, Serrivomeridae, of eels found in temperate and tropical seas worldwide. Sawtooth eels get their name from the saw-like arrangement of inward-slanting teeth attached to the vomer bone in the roof of the mouth. They are deep water pelagic fish. There are eleven species of sawtooth eels in two genera:
| Biology and health sciences | Anguilliformes | Animals |
4392867 | https://en.wikipedia.org/wiki/Phenylpropanoid | Phenylpropanoid | The phenylpropanoids are a diverse family of organic compounds that are biosynthesized by plants from the amino acids phenylalanine and tyrosine in the shikimic acid pathway. Their name is derived from the six-carbon, aromatic phenyl group and the three-carbon propene tail of coumaric acid, which is the central intermediate in phenylpropanoid biosynthesis. From 4-coumaroyl-CoA emanates the biosynthesis of myriad natural products including lignols (precursors to lignin and lignocellulose), flavonoids, isoflavonoids, coumarins, aurones, stilbenes, catechin, and phenylpropanoids. The coumaroyl component is produced from cinnamic acid.
Phenylpropanoids are found throughout the plant kingdom, where they serve as essential components of a number of structural polymers, provide protection from ultraviolet light, defend against herbivores and pathogens, and also mediate plant-pollinator interactions as floral pigments and scent compounds.
Hydroxycinnamic acids
Phenylalanine is first converted to cinnamic acid by the action of the enzyme phenylalanine ammonia-lyase (PAL). Some plants, mainly monocotyledonous, use tyrosine to synthesize p-coumaric acid by the action of the bifunctional enzyme phenylalanine/tyrosine ammonia-lyase (PTAL). A series of enzymatic hydroxylations and methylations leads to coumaric acid, caffeic acid, ferulic acid, 5-hydroxyferulic acid, and sinapic acid. Conversion of these acids to their corresponding esters produces some of the volatile components of herb and flower fragrances, which serve many functions such as attracting pollinators. Ethyl cinnamate is a common example.
Cinnamic aldehydes and monolignols
Reduction of the carboxylic acid functional groups in the cinnamic acids provides the corresponding aldehydes, such as cinnamaldehyde. Further reduction provides monolignols including coumaryl alcohol, coniferyl alcohol, and sinapyl alcohol, which vary only in their degree of methoxylation. The monolignols are monomers that are polymerized to generate various forms of lignin and suberin, which are used as a structural component of plant cell walls.
The phenylpropenes, phenylpropanoids with allylbenzene (3-phenylpropene) as the parent compound, are also derived from the monolignols. Examples include eugenol, chavicol, safrole, and estragole. These compounds are the primary constituents of various essential oils.
Coumarins and flavonoids
Hydroxylation of cinnamic acid in the 4-position by trans-cinnamate 4-monooxygenase leads to p-coumaric acid, which can be further modified into hydroxylated derivatives such as umbelliferone. Another use of p-coumaric acid via its thioester with coenzyme A, i.e. 4-coumaroyl-CoA, is the production of chalcones. This is achieved with the addition of three malonyl-CoA molecules and their cyclization into a second phenyl group. Chalcones are the precursors of all flavonoids, a diverse class of phytochemicals.
Stilbenoids
Stilbenoids, such as resveratrol, are hydroxylated derivatives of stilbene. They are formed through an alternative cyclization of cinnamoyl-CoA or 4-coumaroyl-CoA.
Sporopollenin
Phenylpropanoids and other phenolics are part of the chemical composition of sporopollenin. It is related to cutin and suberin. This ill-defined substance found in pollen is unusually resistant to degradation. Analyses have revealed a mixture of biopolymers, containing mainly hydroxylated fatty acids, phenylpropanoids, phenolics and traces of carotenoids. Tracer experiments have shown that phenylalanine is a major precursor, but other carbon sources also contribute. It is likely that sporopollenin is derived from several precursors that are chemically cross-linked to form a rigid structure.
| Physical sciences | Phenylpropanoids | Chemistry |
4396171 | https://en.wikipedia.org/wiki/Earth%27s%20rotation | Earth's rotation | Earth's rotation or Earth's spin is the rotation of planet Earth around its own axis, as well as changes in the orientation of the rotation axis in space. Earth rotates eastward, in prograde motion. As viewed from the northern polar star Polaris, Earth turns counterclockwise.
The North Pole, also known as the Geographic North Pole or Terrestrial North Pole, is the point in the Northern Hemisphere where Earth's axis of rotation meets its surface. This point is distinct from Earth's north magnetic pole. The South Pole is the other point where Earth's axis of rotation intersects its surface, in Antarctica.
Earth rotates once in about 24 hours with respect to the Sun, but once every 23 hours, 56 minutes and 4 seconds with respect to other distant stars (see below). Earth's rotation is slowing slightly with time; thus, a day was shorter in the past. This is due to the tidal effects the Moon has on Earth's rotation. Atomic clocks show that the modern day is longer by about 1.7 milliseconds than a century ago, slowly increasing the rate at which UTC is adjusted by leap seconds. Analysis of historical astronomical records shows a slowing trend; the length of a day increased by about 2.3 milliseconds per century since the 8th century BCE.
Scientists reported that in 2020 Earth had started spinning faster, after consistently spinning slower than 86,400 seconds per day in the decades before. On June 29, 2022, Earth's spin was completed in 1.59 milliseconds under 24 hours, setting a new record. Because of that trend, engineers worldwide are discussing a 'negative leap second' and other possible timekeeping measures.
This increase in speed is thought to be due to various factors, including the complex motion of its molten core, oceans, and atmosphere, the effect of celestial bodies such as the Moon, and possibly climate change, which is causing the ice at Earth's poles to melt. The masses of ice account for the Earth's shape being that of an oblate spheroid, bulging around the equator. When these masses are reduced, the poles rebound from the loss of weight, and Earth becomes more spherical, which has the effect of bringing mass closer to its centre of gravity. Conservation of angular momentum dictates that a mass distributed more closely around its centre of gravity spins faster.
History
Among the ancient Greeks, several of the Pythagorean school believed in the rotation of Earth rather than the apparent diurnal rotation of the heavens. Perhaps the first was Philolaus (470–385 BCE), though his system was complicated, including a counter-earth rotating daily about a central fire.
A more conventional picture was supported by Hicetas, Heraclides and Ecphantus in the fourth century BCE who assumed that Earth rotated but did not suggest that Earth revolved about the Sun. In the third century BCE, Aristarchus of Samos suggested the Sun's central place.
However, Aristotle in the fourth century BCE criticized the ideas of Philolaus as being based on theory rather than observation. He established the idea of a sphere of fixed stars that rotated about Earth. This was accepted by most of those who came after, in particular Claudius Ptolemy (2nd century CE), who thought Earth would be devastated by gales if it rotated.
In 499 CE, the Indian astronomer Aryabhata suggested that the spherical Earth rotates about its axis daily and that the apparent movement of the stars is a relative motion caused by the rotation of the Earth. He provided the following analogy: "Just as a man in a boat going in one direction sees the stationary things on the bank as moving in the opposite direction, in the same way to a man at Lanka the fixed stars appear to be going westward."
In the 10th century, some Muslim astronomers accepted that the Earth rotates around its axis. According to al-Biruni, al-Sijzi (d. c. 1020) invented an astrolabe called al-zūraqī based on the idea believed by some of his contemporaries "that the motion we see is due to the Earth's movement and not to that of the sky." The prevalence of this view is further confirmed by a reference from the 13th century which states: "According to the geometers [or engineers] (muhandisīn), the Earth is in constant circular motion, and what appears to be the motion of the heavens is actually due to the motion of the Earth and not the stars." Treatises were written to discuss its possibility, either as refutations or expressing doubts about Ptolemy's arguments against it. At the Maragha and Samarkand observatories, Earth's rotation was discussed by Tusi (born 1201) and Qushji (born 1403); the arguments and evidence they used resemble those used by Copernicus.
In medieval Europe, Thomas Aquinas accepted Aristotle's view and so, reluctantly, did John Buridan and Nicole Oresme in the fourteenth century. Not until Nicolaus Copernicus in 1543 adopted a heliocentric world system did the contemporary understanding of Earth's rotation begin to be established. Copernicus pointed out that if the movement of the Earth is violent, then the stars' movement must be much more so. He acknowledged the contribution of the Pythagoreans and pointed to examples of relative motion. For Copernicus, this was the first step in establishing the simpler pattern of planets circling a central Sun.
Tycho Brahe, who produced accurate observations on which Kepler based his laws of planetary motion, used Copernicus's work as the basis of a system assuming a stationary Earth. In 1600, William Gilbert strongly supported Earth's rotation in his treatise on Earth's magnetism and thereby influenced many of his contemporaries. Those like Gilbert who did not openly support or reject the motion of Earth about the Sun are called "semi-Copernicans". A century after Copernicus, Riccioli disputed the model of a rotating Earth due to the lack of then-observable eastward deflections in falling bodies; such deflections would later be called the Coriolis effect. However, the contributions of Kepler, Galileo, and Newton gathered support for the theory of the rotation of the Earth.
Empirical tests
Earth's rotation implies that the Equator bulges and the geographical poles are flattened.
In his Principia, Newton predicted this flattening would amount to one part in 230, and pointed to the pendulum measurements taken by Richer in 1673 as corroboration of the change in gravity, but initial measurements of meridian lengths by Picard and Cassini at the end of the 17th century suggested the opposite.
However, measurements by Maupertuis and the French Geodesic Mission in the 1730s established the oblateness of Earth, thus confirming the positions of both Newton and Copernicus.
In Earth's rotating frame of reference, a freely moving body follows an apparent path that deviates from the one it would follow in a fixed frame of reference. Because of the Coriolis effect, falling bodies veer slightly eastward from the vertical plumb line below their point of release, and projectiles veer right in the Northern Hemisphere (and left in the Southern) from the direction in which they are shot. The Coriolis effect is mainly observable at a meteorological scale, where it is responsible for the opposite directions of cyclone rotation in the Northern and Southern hemispheres (anticlockwise and clockwise, respectively).
Hooke, following a suggestion from Newton in 1679, tried unsuccessfully to verify the predicted eastward deviation of a body dropped from a height of , but definitive results were obtained later, in the late 18th and early 19th centuries, by Giovanni Battista Guglielmini in Bologna, Johann Friedrich Benzenberg in Hamburg and Ferdinand Reich in Freiberg, using taller towers and carefully released weights. A ball dropped from a height of 158.5 m departed by 27.4 mm from the vertical compared with a calculated value of 28.1 mm.
The most celebrated test of Earth's rotation is the Foucault pendulum first built by physicist Léon Foucault in 1851, which consisted of a lead-filled brass sphere suspended from the top of the Panthéon in Paris. Because of Earth's rotation under the swinging pendulum, the pendulum's plane of oscillation appears to rotate at a rate depending on latitude. At the latitude of Paris, the predicted and observed shift was about clockwise per hour. Foucault pendulums now swing in museums worldwide.
Periods
True solar day
Earth's rotation period relative to the Sun (solar noon to solar noon) is its true solar day or apparent solar day. It depends on Earth's orbital motion and is thus affected by changes in the eccentricity and inclination of Earth's orbit. Both vary over thousands of years, so the annual variation of the true solar day also varies. Generally, it is longer than the mean solar day during two periods of the year and shorter during another two. The true solar day tends to be longer near perihelion when the Sun apparently moves along the ecliptic through a greater angle than usual, taking about longer to do so. Conversely, it is about shorter near aphelion. It is about longer near a solstice when the projection of the Sun's apparent motion along the ecliptic onto the celestial equator causes the Sun to move through a greater angle than usual. Conversely, near an equinox the projection onto the equator is shorter by about . Currently, the perihelion and solstice effects combine to lengthen the true solar day near by solar seconds, but the solstice effect is partially cancelled by the aphelion effect near when it is only longer. The effects of the equinoxes shorten it near and by and , respectively.
Mean solar day
The average of the true solar day during the course of an entire year is the mean solar day, which contains 86,400 mean solar seconds. Currently, each of these seconds is slightly longer than an SI second because Earth's mean solar day is now slightly longer than it was during the 19th century due to tidal friction. The average length of the mean solar day since the introduction of the leap second in 1972 has been about 0 to 2 ms longer than 86,400 SI seconds. Random fluctuations due to core-mantle coupling have an amplitude of about 5 ms. The mean solar second between 1750 and 1892 was chosen in 1895 by Simon Newcomb as the independent unit of time in his Tables of the Sun. These tables were used to calculate the world's ephemerides between 1900 and 1983, so this second became known as the ephemeris second. In 1967 the SI second was made equal to the ephemeris second.
The apparent solar time is a measure of Earth's rotation and the difference between it and the mean solar time is known as the equation of time.
Stellar and sidereal day
Earth's rotation period relative to the International Celestial Reference Frame, called its stellar day by the International Earth Rotation and Reference Systems Service (IERS), is seconds of mean solar time (UT1) , ). Earth's rotation period relative to the precessing mean vernal equinox, named sidereal day, is of mean solar time (UT1) , ). Thus, the sidereal day is shorter than the stellar day by about .
Both the stellar day and the sidereal day are shorter than the mean solar day by about . This is a result of the Earth turning 1 additional rotation, relative to the celestial reference frame, as it orbits the Sun (so 366.24 rotations/y). The mean solar day in SI seconds is available from the IERS for the periods and .
Recently (1999–2010) the average annual length of the mean solar day in excess of 86,400 SI seconds has varied between and , which must be added to both the stellar and sidereal days given in mean solar time above to obtain their lengths in SI seconds (see Fluctuations in the length of day).
Angular speed
The angular speed of Earth's rotation in inertial space is ± . Multiplying by (180°/π radians) × (86,400 seconds/day) yields , indicating that Earth rotates more than 360 degrees relative to the fixed stars in one solar day. Earth's movement along its nearly circular orbit while it is rotating once around its axis requires that Earth rotate slightly more than once relative to the fixed stars before the mean Sun can pass overhead again, even though it rotates only once (360°) relative to the mean Sun. Multiplying the value in rad/s by Earth's equatorial radius of (WGS84 ellipsoid) (factors of 2π radians needed by both cancel) yields an equatorial speed of . Some sources state that Earth's equatorial speed is slightly less, or . This is obtained by dividing Earth's equatorial circumference by . However, the use of the solar day is incorrect; it must be the sidereal day, so the corresponding time unit must be a sidereal hour. This is confirmed by multiplying by the number of sidereal days in one mean solar day, , which yields the equatorial speed in mean solar hours given above of 1,674.4 km/h.
The tangential speed of Earth's rotation at a point on Earth can be approximated by multiplying the speed at the equator by the cosine of the latitude. For example, the Kennedy Space Center is located at latitude 28.59° N, which yields a speed of: cos(28.59°) × 1,674.4 km/h = 1,470.2 km/h. Latitude is a placement consideration for spaceports.
The peak of the Cayambe volcano is the point of Earth's surface farthest from its axis; thus, it rotates the fastest as Earth spins.
Changes
In rotational axis
Earth's rotation axis moves with respect to the fixed stars (inertial space); the components of this motion are precession and nutation. It also moves with respect to Earth's crust; this is called polar motion.
Precession is a rotation of Earth's rotation axis, caused primarily by external torques from the gravity of the Sun, Moon and other bodies. The polar motion is primarily due to free core nutation and the Chandler wobble.
In rotational speed
Tidal interactions
Over millions of years, Earth's rotation has been slowed significantly by tidal acceleration through gravitational interactions with the Moon. Thus angular momentum is slowly transferred to the Moon at a rate proportional to , where is the orbital radius of the Moon. This process has gradually increased the length of the day to its current value, and resulted in the Moon being tidally locked with Earth.
This gradual rotational deceleration is empirically documented by estimates of day lengths obtained from observations of tidal rhythmites and stromatolites; a compilation of these measurements found that the length of the day has increased steadily from about 21 hours at 600 Myr ago to the current 24-hour value. By counting the microscopic lamina that form at higher tides, tidal frequencies (and thus day lengths) can be estimated, much like counting tree rings, though these estimates can be increasingly unreliable at older ages.
Resonant stabilization
The current rate of tidal deceleration is anomalously high, implying Earth's rotational velocity must have decreased more slowly in the past. Empirical data tentatively shows a sharp increase in rotational deceleration about 600 Myr ago. Some models suggest that Earth maintained a constant day length of 21 hours throughout much of the Precambrian. This day length corresponds to the semidiurnal resonant period of the thermally driven atmospheric tide; at this day length, the decelerative lunar torque could have been canceled by an accelerative torque from the atmospheric tide, resulting in no net torque and a constant rotational period. This stabilizing effect could have been broken by a sudden change in global temperature. Recent computational simulations support this hypothesis and suggest the Marinoan or Sturtian glaciations broke this stable configuration about 600 Myr ago; the simulated results agree quite closely with existing paleorotational data.
Global events
Some recent large-scale events, such as the 2004 Indian Ocean earthquake, have caused the length of a day to shorten by 3 microseconds by reducing Earth's moment of inertia. Post-glacial rebound, ongoing since the last ice age, is also changing the distribution of Earth's mass, thus affecting the moment of inertia of Earth and, by the conservation of angular momentum, Earth's rotation period.
The length of the day can also be influenced by man-made structures. For example, NASA scientists calculated that the water stored in the Three Gorges Dam has increased the length of Earth's day by 0.06 microseconds due to the shift in mass.
Measurement
The primary monitoring of Earth's rotation is performed by very-long-baseline interferometry coordinated with the Global Positioning System, satellite laser ranging, and other satellite geodesy techniques. This provides an absolute reference for the determination of universal time, precession and nutation.
The absolute value of Earth rotation including UT1 and nutation can be determined using space geodetic observations, such as very-long-baseline interferometry and lunar laser ranging, whereas their derivatives, denoted as length-of-day excess and nutation rates can be derived from satellite observations, such as GPS, GLONASS, Galileo and satellite laser ranging to geodetic satellites.
Ancient observations
There are recorded observations of solar and lunar eclipses by Babylonian and Chinese astronomers beginning in the 8th century BCE, as well as from the medieval Islamic world and elsewhere. These observations can be used to determine changes in Earth's rotation over the last 27 centuries, since the length of the day is a critical parameter in the calculation of the place and time of eclipses. A change in day length of milliseconds per century shows up as a change of hours and thousands of kilometers in eclipse observations. The ancient data are consistent with a shorter day, meaning Earth was turning faster throughout the past.
Cyclic variability
Around every 25–30 years Earth's rotation slows temporarily by a few milliseconds per day, usually lasting around five years. 2017 was the fourth consecutive year that Earth's rotation has slowed. The cause of this variability has not yet been determined.
Origin
Earth's original rotation was a vestige of the original angular momentum of the cloud of dust, rocks and gas that coalesced to form the Solar System. This primordial cloud was composed of hydrogen and helium produced in the Big Bang, as well as heavier elements ejected by supernovas. As this interstellar dust is heterogeneous, any asymmetry during gravitational accretion resulted in the angular momentum of the eventual planet.
However, if the giant-impact hypothesis for the origin of the Moon is correct, this primordial rotation rate would have been reset by the Theia impact 4.5 billion years ago. Regardless of the speed and tilt of Earth's rotation before the impact, it would have experienced a day some five hours long after the impact. Tidal effects would then have slowed this rate to its modern value.
| Physical sciences | Earth science basics: General | Earth science |
4396843 | https://en.wikipedia.org/wiki/Wood%20drying | Wood drying | Wood drying (also seasoning lumber or wood seasoning) reduces the moisture content of wood before its use. When the drying is done in a kiln, the product is known as kiln-dried timber or lumber, whereas air drying is the more traditional method.
There are two main reasons for drying wood:
WoodworkingWhen wood is used as a construction material, whether as a structural support in a building or in woodworking objects, it will absorb or expel moisture until it is in equilibrium with its surroundings. Equilibration (usually drying) causes unequal shrinkage in the wood, and can cause damage to the wood if equilibration occurs too rapidly. The equilibration must be controlled to prevent damage to the wood.
Wood burning When wood is burned (firewood), it is usually best to dry it first. Damage from shrinkage is not a problem here, as it may be in the case of drying for woodworking purposes. Moisture affects the burning process, with unburnt hydrocarbons going up the chimney. If a 50% wet log is burnt at high temperature, with good heat extraction from the exhaust gas leading to a 100 °C exhaust temperature, about 5% of the energy of the log is wasted through evaporating and heating the water vapour. With condensers, the efficiency can be further increased; but, for the normal stove, the key to burning wet wood is to burn it very hot, perhaps starting fire with dry wood.
For some purposes, wood is not dried at all, and is used green. Often, wood must be in equilibrium with the air outside, as for construction wood, or the air indoors, as for wooden furniture.
Wood is air-dried or dried in a purpose built oven (kiln). Usually the wood is sawn before drying, but sometimes the log is dried whole.
Case hardening describes lumber or timber that has been dried too rapidly. Wood initially dries from the shell (surface), shrinking the shell and putting the core under compression. When this shell is at a low moisture content it will 'set' and resist shrinkage. The core of the wood is still at a higher moisture content. This core will then begin to dry and shrink. However, any shrinkage is resisted by the already 'set' shell. This leads to reversed stresses; compression stresses on the shell and tension stresses in the core. This results in unrelieved stress called case hardening. Case-hardened [wood] may warp considerably and dangerously when the stress is released by sawing.
Types of wood
Wood is divided, according to its botanical origin, into two kinds: softwoods, from coniferous trees, and hardwoods, from broad-leaved trees. Softwoods are lighter and generally simple in structure, whereas hardwoods are harder and more complex. However, in Australia, softwood generally describes rain forest trees, and hardwood describes Sclerophyll species (Eucalyptus spp).
Softwoods such as pine are typically much lighter and easier to process than hardwoods such as fruit tree wood. The density of softwoods ranges from to , while hardwoods are to . Once dried, both consist of approximately 12% of moisture (Desch and Dinwoodie, 1996). Because of hardwood's denser and more complex structure, its permeability is much less than that of softwood, making it more difficult to dry. Although there are about a hundred times more species of hardwood trees than softwood trees, the ability to be dried and processed faster and more easily makes softwood the main supply of commercial wood nowadays.
Wood–water relationships
The timber of living trees and fresh logs contains a large amount of water which often constitutes over 50% of the wood's weight. Water has a significant influence on wood. Wood continually exchanges moisture or water with its surroundings, although the rate of exchange is strongly affected by the degree to which wood is sealed.
Wood contains water in three forms:
Free water The bulk of water contained in the cell lumina is only held by capillary forces. It is not bound chemically and is called free water. Free water is not in the same thermodynamic state as liquid water: energy is required to overcome the capillary forces. Furthermore, free water may contain chemicals, altering the drying characteristics of wood.
Bound or hygroscopic water Bound water is bound to the wood via hydrogen bonds. The attraction of wood for water arises from the presence of free hydroxyl (OH) groups in the cellulose, hemicelluloses and lignin molecules in the cell wall. The hydroxyl groups are negatively charged. Because water is a polar liquid, the free hydroxyl groups in cellulose attract and hold water by hydrogen bonding.
Vapor Water in cell lumina in the form of water vapour is normally negligible at normal temperature and humidity.
Moisture content
The moisture content of wood is calculated as the mass change as a proportion of the dry mass, by the formula (Siau, 1984):
Here, is the green mass of the wood, is its oven dry mass (the attainment of constant mass generally after drying in an oven set at () for 24 hours as mentioned by Walker et al., 1993). The equation can also be expressed as a fraction of the mass of the water and the mass of the oven dry wood rather than a percentage. For example, (oven dry basis) expresses the same moisture content as 59% (oven dry basis).
Fibre saturation point
When green wood dries, free water from the cell lumina, held by the capillary forces only, is the first to go. Physical properties, such as strength and shrinkage, are generally not affected by the removal of free water. The fibre saturation point (FSP) is defined as the moisture content at which free water should be completely gone, while the cell walls are saturated with bound water. In most types of woods, the fibre saturation point is at 25 to 30% moisture content. Siau (1984) reported that the fibre saturation point (kg/kg) is dependent on the temperature T (°C) according to the following equation:
(1.2)
Keey et al. (2000) use a different definition of the fibre saturation point (equilibrium moisture content of wood in an environment of 99% relative humidity).
Many properties of wood show considerable change as the wood is dried below the fibre saturation point, including:
volume (ideally no shrinkage occurs until some bound water is lost, that is, until wood is dried below FSP);
strength (strengths generally increase consistently as the wood is dried below the FSP (Desch and Dinwoodie, 1996), except for impact-bending strength and, in some cases, toughness);
electrical resistivity, which increases very rapidly with the loss of bound water when the wood dries below the FSP.
Equilibrium moisture content
Wood is a hygroscopic substance. It has the ability to take in or give off moisture in the form of vapour. Water contained in wood exerts vapour pressure of its own, which is determined by the maximum size of the capillaries filled with water at any time. If water vapour pressure in the ambient space is lower than vapour pressure within wood, desorption takes place. The largest-sized capillaries, which are full of water at the time, empty first. Vapour pressure within the wood falls as water is successively contained in smaller capillaries. A stage is eventually reached when vapour pressure within the wood equals vapour pressure in the ambient space above the wood, and further desorption ceases. The amount of moisture that remains in the wood at this stage is in equilibrium with water vapour pressure in the ambient space, and is termed the equilibrium moisture content or EMC (Siau, 1984). Because of its hygroscopicity, wood tends to reach a moisture content that is in equilibrium with the relative humidity and temperature of the surrounding air.
The EMC of wood varies with the ambient relative humidity (a function of temperature) significantly, to a lesser degree with the temperature. Siau (1984) reported that the EMC also varies very slightly with species, mechanical stress, drying history of wood, density, extractives content and the direction of sorption in which the moisture change takes place (i.e. adsorption or desorption).
Moisture content of wood in service
Wood retains its hygroscopic characteristics after it is put into use. It is then subjected to fluctuating humidity, the dominant factor in determining its EMC. These fluctuations may be more or less cyclical, such as diurnal changes or annual seasonal changes.
To minimize the changes in wood moisture content or the movement of wooden objects in service, wood is usually dried to a moisture content that is close to the average EMC conditions to which it will be exposed. These conditions vary for interior uses compared with exterior uses in a given geographic location. For example, according to the Australian Standard for Timber Drying Quality (AS/NZS 4787, 2001), the EMC is recommended to be 10–12% for the majority of Australian states, although extreme cases are up to 15 to 18% for some places in Queensland, Northern Territory, Western Australia and Tasmania. However, the EMC is as low as 6 to 7% in dry centrally heated houses and offices or in permanently air-conditioned buildings.
Shrinkage and swelling
Shrinkage and swelling may occur in wood when the moisture content is changed (Stamm, 1964). Shrinkage occurs as moisture content decreases, while swelling takes place when it increases. Volume change is not equal in all directions. The greatest dimensional change occurs in a direction tangential to the growth rings. Shrinkage from the pith outwards, or radially, is usually considerably less than tangential shrinkage, while longitudinal (along the grain) shrinkage is so slight as to be usually neglected. The longitudinal shrinkage is 0.1% to 0.3%, in contrast to transverse shrinkages, which is 2% to 10%. Tangential shrinkage is often about twice as great as in the radial direction, although in some species it is as much as five times as great. The shrinkage is about 5% to 10% in the tangential direction and about 2% to 6% in the radial direction (Walker et al., 1993).
Differential transverse shrinkage of wood is related to:
the alternation of late wood and early wood increments within the annual ring;
the influence of wood rays on the radial direction (Kollmann and Cote, 1968);
the features of the cell wall structure such as microfibril angle modifications and pits;
the chemical composition of the middle lamella.
Wood drying may be described as the art of ensuring that gross dimensional changes through shrinkage are confined to the drying process. Ideally, wood is dried to that equilibrium moisture content as will later (in service) be attained by the wood. Thus, further dimensional change will be kept to a minimum.
It is probably impossible to completely eliminate dimensional change in wood, but elimination of change in size may be approximated by chemical modification. For example, wood can be treated with chemicals to replace the hydroxyl groups with other hydrophobic functional groups of modifying agents (Stamm, 1964). Among all the existing processes, wood modification with acetic anhydride has been noted for the high anti-shrink or anti-swell efficiency (ASE) attainable without damage to wood. However, acetylation of wood has been slow to be commercialised due to the cost, corrosion and the entrapment of the acetic acid in wood. There is an extensive volume of literature relating to the chemical modification of wood (Rowell, 1983, 1991; Kumar, 1994; Haque, 1997).
Drying timber is one method of adding value to sawn products from the primary wood processing industries. According to the Australian Forest and Wood Products Research and Development Corporation (FWPRDC), green sawn hardwood, which is sold at about $350 per cubic metre or less, increases in value to $2,000 per cubic metre or more with drying and processing. However, currently used conventional drying processes often result in significant quality problems from cracks, both externally and internally, reducing the value of the product. For example, in Queensland (Anon, 1997), on the assumption that 10% of the dried softwood is devalued by $200 per cubic metre because of drying defects, saw millers are losing about $5 million a year. In Australia, the loss could be $40 million a year for softwood and an equal or higher amount for hardwood. Thus, proper drying under controlled conditions prior to use is of great importance in timber use, in countries where climatic conditions vary considerably at different times of the year.
Drying, if carried out promptly after felling of trees, also protects timber against primary decay, fungal stain and attack by certain kinds of insects. Organisms, which cause decay and stain, generally cannot thrive in timber with a moisture content below 20%. Several, though not all, insect pests can live only in green timber.
In addition to the above advantages of drying timber, the following points are also significant (Walker et al., 1993; Desch and Dinwoodie, 1996):
Dried timber is lighter, and the transportation and handling costs are reduced.
Dried timber is stronger than green timber in most strength properties.
Timbers for impregnation with preservatives have to be properly dried if proper penetration is to be accomplished, particularly in the case of oil-type preservatives.
In the field of chemical modification of wood and wood products, the material should be dried to a certain moisture content for the appropriate reactions to occur.
Dry wood generally works, machines, finishes and glues better than green timber (although there are exceptions; for instance, green wood is often easier to turn than dry wood). Paints and finishes last longer on dry timber.
The electrical and thermal insulation properties of wood are improved by drying.
Prompt drying of wood immediately after felling therefore significantly upgrades and adds value to raw timber. Drying enables substantial long-term economy by rationalizing the use of timber resources. The drying of wood is thus an area for research and development, which concern many researchers and timber companies around the world.
Mechanisms of moisture movement
Water in wood normally moves from zones of higher to zones of lower moisture content (Walker et al., 1993). Drying starts from the exterior of the wood and moves towards the centre, and drying at the outside is also necessary to expel moisture from the inner zones of the wood. Wood subsequently attains equilibrium with the surrounding air in moisture content.
Moisture passageways
The driving force of moisture movement is chemical potential. However, it is not always easy to relate chemical potential in wood to commonly observable variables, such as temperature and moisture content (Keey et al., 2000). Moisture in wood moves within the wood as liquid or vapour through several types of passageways, based on the nature of the driving force, (e.g. pressure or moisture gradient), and variations in wood structure (Langrish and Walker, 1993), as explained in the next section on driving forces for moisture movement. These pathways consist of cavities of the vessels, fibres, ray cells, pit chambers and their pit membrane openings, intercellular spaces and transitory cell wall passageways.
Movement of water takes place in these passageways in any direction, longitudinally in the cells, as well as laterally from cell to cell until it reaches the lateral drying surfaces of the wood. The higher longitudinal permeability of sapwood of hardwood is generally caused by the presence of vessels. The lateral permeability and transverse flow is often very low in hardwoods. The vessels in hardwoods are sometimes blocked by the presence of tyloses and/or by secreting gums and resins in some other species, as mentioned earlier. The presence of gum veins, the formation of which is often a result of natural protective response of trees to injury, is commonly observed on the surface of sawn boards of most eucalypts. Despite the generally higher volume fraction of rays in hardwoods (typically 15% of wood volume), the rays are not particularly effective in radial flow, nor are the pits on the radial surfaces of fibres effective in tangential flow (Langrish and Walker, 1993).
Moisture movement space
The available space for air and moisture in wood depends on the density and porosity of wood. Porosity is the volume fraction of void space in a solid. The porosity is reported to be 1.2 to 4.6% of dry volume of wood cell wall (Siau, 1984). On the other hand, permeability is a measure of the ease with which fluids are transported through a porous solid under the influence of some driving forces, e.g. capillary pressure gradient or moisture gradient. It is clear that solids must be porous to be permeable, but it does not necessarily follow that all porous bodies are permeable. Permeability can only exist if the void spaces are interconnected by openings. For example, a hardwood may be permeable because there is intervessel pitting with openings in the membranes (Keey et al., 2000). If these membranes are occluded or encrusted, or if the pits are aspirated, the wood assumes a closed-cell structure and may be virtually impermeable. The density is also important for impermeable hardwoods because more cell-wall material is traversed per unit distance, which offers increased resistance to diffusion (Keey et al., 2000). Hence lighter woods, in general, dry more rapidly than do the heavier woods. The transport of fluids is often bulk flow (momentum transfer) for permeable softwoods at high temperature while diffusion occurs for impermeable hardwoods (Siau, 1984). These mechanisms are discussed below.
Driving forces for moisture movement
Three main driving forces used in different version of diffusion models are moisture content, the partial pressure of water vapour, and the chemical potential of water (Skaar, 1988; Keey et al., 2000). These are discussed here, including capillary action, which is a mechanism for free water transport in permeable softwoods. Total pressure difference is the driving force during wood vacuum drying.
Capillary action
Capillary forces determine the movements (or absence of movement) of free water. It is due to both adhesion and cohesion. Adhesion is the attraction between water to other substances and cohesion is the attraction of the molecules in water to each other.
As wood dries, evaporation of water from the surface sets up capillary forces that exert a pull on the free water in the zones of wood beneath the surfaces. When there is no longer any free water in the wood capillary forces are no longer of importance.
Moisture content differences
The chemical potential is explained here since it is the true driving force for the transport of water in both liquid and vapour phases in wood (Siau, 1984). The Gibbs free energy per mole of substance is usually expressed as the chemical potential of that substance (Skaar, 1933). The chemical potential of water in unsaturated air or wood below the fibre saturation point influences the drying of wood. Equilibrium will occur at the equilibrium moisture content (as defined earlier) of wood when the chemical potential of water in the wood becomes equal to that in the surrounding air. The chemical potential of sorbed water is a function of wood moisture content. Therefore, a gradient of wood moisture content (between surface and centre), or more specifically of water activity, is accompanied by a gradient of chemical potential under isothermal conditions. Moisture will redistribute itself throughout the wood until its chemical potential is uniform throughout, resulting in a zero potential gradient at equilibrium (Skaar, 1988). The flux of moisture attempting to achieve the equilibrium state is assumed to be proportional to the difference in its chemical potential, and inversely proportional to the path length over which the potential difference acts (Keey et al., 2000).
The gradient in chemical potential is related to the moisture content gradient as explained in above equations (Keey et al., 2000). The diffusion model using the moisture content gradient as a driving force was applied successfully by Wu (1989) and Doe et al. (1994). Though the agreement between the moisture-content profiles predicted by the diffusion model based on moisture-content gradients is better at lower moisture contents than at higher ones, there is no evidence to suggest that there are significantly different moisture-transport mechanisms operating at higher moisture contents for this timber. Their observations are consistent with a transport process that is driven by the total concentration of water. The diffusion model is used here based on this empirical evidence that the moisture-content gradient is a driving force for drying this type of impermeable timber.
Differences in moisture content between the surface and the centre (gradient, the chemical potential difference between interface and bulk) move the bound water through the small passageways in the cell wall by diffusion. In comparison with capillary movement, diffusion is a slow process. Diffusion is the generally suggested mechanism for the drying of impermeable hardwoods (Keey et al., 2000). Furthermore, moisture migrates slowly due to the fact that extractives plug the small cell wall openings in the heartwood. This is why sapwood generally dries faster than heartwood under the same drying conditions.
Moisture movement directions for diffusion
It is reported that the ratio of the longitudinal to the transverse (radial and tangential) diffusion rates for wood ranges from about 100 at a moisture content of 5%, to 2–4 at a moisture content of 25% (Langrish and Walker, 1993). Radial diffusion is somewhat faster than tangential diffusion. Although longitudinal diffusion is most rapid, it is of practical importance only when short pieces are dried. Generally the timber boards are much longer than in width or thickness. For example, a typical size of a green board used for this research was 6m long, 250 mm in width and 43 mm in thickness. If the boards are quartersawn, then the width will be in the radial direction whereas the thickness will be in tangential direction, and vice versa for plain-sawn boards. Most of the moisture is removed from wood by lateral movement during drying.
Reasons for splits and cracks during timber drying and their control
The chief difficulty experienced in the drying of timber is the tendency of its outer layers to dry out more rapidly than the interior ones. If these layers are allowed to dry much below the fibre saturation point while the interior is still saturated, stresses (called drying stresses) are set up because the shrinkage of the outer layers is restricted by the wet interior (Keey et al., 2000). Rupture in the wood tissues occurs, and consequently splits and cracks occur if these stresses across the grain exceed the strength across the grain (fibre to fibre bonding).
The successful control of drying defects in a drying process consists in maintaining a balance between the rate of evaporation of moisture from the surface and the rate of outward movement of moisture from the interior of the wood. The way in which drying can be controlled will now be explained. One of the most successful ways of wood drying or seasoning would be kiln drying, where the wood is placed into a kiln compartment in stacks and dried by steaming, and releasing the steam slowly.
Influence of temperature, relative humidity and rate of air circulation
The external drying conditions (temperature, relative humidity and air velocity) control the external boundary conditions for drying, and hence the drying rate, as well as affecting the rate of internal moisture movement. The drying rate is affected by external drying conditions (Walker et al., 1993; Keey et al., 2000), as will now be described.
Temperature If the relative humidity is kept constant, the higher the temperature, the higher the drying rate. Temperature influences the drying rate by increasing the moisture holding capacity of the air, as well as by accelerating the diffusion rate of moisture through the wood. The actual temperature in a drying kiln is the dry-bulb temperature (usually denoted by Tg), which is the temperature of a vapour-gas mixture determined by inserting a thermometer with a dry bulb. On the other hand, the wet-bulb temperature (TW) is defined as the temperature reached by a small amount of liquid evaporating in a large amount of an unsaturated air-vapour mixture. The temperature sensing element of this thermometer is kept moist with a porous fabric sleeve (cloth) usually put in a reservoir of clean water. A minimum air flow of 2 m/s is needed to prevent a zone of stagnant damp air formation around the sleeve (Walker et al., 1993). Since air passes over the wet sleeve, water is evaporated and cools the wet-bulb thermometer. The difference between the dry-bulb and wet-bulb temperatures, the wet-bulb depression, is used to determine the relative humidity from a standard hygrometric chart (Walker et al., 1993). A higher difference between the dry-bulb and wet-bulb temperatures indicates a lower relative humidity. For example, if the dry-bulb temperature is 100 °C and wet-bulb temperature 60 °C, then the relative humidity is read as 17% from a hygrometric chart.
Relative humidity The relative humidity of air is defined as the partial pressure of water vapour divided by the saturated vapour pressure at the same temperature and total pressure (Siau, 1984). If the temperature is kept constant, lower relative humidities result in higher drying rates due to the increased moisture gradient in wood, resulting from the reduction of the moisture content in the surface layers when the relative humidity of air is reduced. The relative humidity is usually expressed on a percentage basis. For drying, the other essential parameter related to relative humidity is the absolute humidity, which is the mass of water vapour per unit mass of dry air (kg of water per kg of dry air). However, its influenced by the amount of water in the heated air.
Air circulation rate Drying time and timber quality depend on the air velocity and its uniform circulation. At a constant temperature and relative humidity, the highest possible drying rate is obtained by rapid circulation of air across the surface of wood, giving rapid removal of moisture evaporating from the wood. However, a higher drying rate is not always desirable, particularly for impermeable hardwoods, because higher drying rates develop greater stresses that may cause the timber to crack or distort. At very low fan speeds, less than 1 m/s, the air flow through the stack is often laminar flow, and the heat transfer between the timber surface and the moving air stream is not particularly effective (Walker et al., 1993). The low effectiveness (externally) of heat transfer is not necessarily a problem if internal moisture movement is the key limitation to the movement of moisture, as it is for most hardwoods (Pordage and Langrish, 1999).
Classification of timbers for drying
The timbers are classified as follows according to their ease of drying and their proneness to drying degrade:
Highly refractory woods These woods are slow and difficult to dry if the final product is to be free from defects, particularly cracks and splits. Examples are heavy structural timbers with high density such as ironbark (Eucalyptus paniculata), blackbutt (E. pillularis), southern blue gum (E. globulus) and brush box (Lophostemon cofertus). They require considerable protection and care against rapid drying conditions for the best results (Bootle, 1994).
Moderately refractory woods These timbers show a moderate tendency to crack and split during seasoning. They can be seasoned free from defects with moderately rapid drying conditions (i.e. a maximum dry-bulb temperature of 85 °C can be used). Examples are Sydney blue gum (E. saligna) and other timbers of medium density (Bootle, 1994), which are potentially suitable for furniture.
Non-refractory woods These woods can be rapidly seasoned to be free from defects even by applying high temperatures (dry-bulb temperatures of more than 100 °C) in industrial kilns. If not dried rapidly, they may develop discolouration (blue stain) and mould on the surface. Examples are softwoods and low density timbers such as Pinus radiata.
Model
The rate at which wood dries depends upon a number of factors, the most important of which are the temperature, the dimensions of the wood, and the relative humidity. Simpson and Tschernitz have developed a simple model of wood drying as a function of these three variables. Although the analysis was done for red oak, the procedure may be applied to any species of wood by adjusting the constant parameters of the model.
Simply put, the model assumes that the rate of change of the moisture content M with respect to time t is proportional to how far the wood sample is from its equilibrium moisture content , which is a function of the temperature T and relative humidity h:
where is a function of the temperature T and a typical wood dimension L and has units of time. The typical wood dimension is roughly the smallest value of () which are the radial, tangential and longitudinal dimensions respectively, in inches, with the longitudinal dimension divided by ten because water diffuses about 10 times more rapidly in the longitudinal direction (along the grain) than in the lateral dimensions. The solution to the above equation is:
Where is the initial moisture content. It was found that for red oak lumber, the "time constant" was well expressed as:
where a, b and n are constants and is the saturation vapor pressure of water at temperature T. For time measured in days, length in inches, and measured in mmHg, the following values of the constants were found for red oak lumber.
a = 0.0575
b = 0.00142
n = 1.52
Solving for the drying time yields:
For example, at 150°F, using the Arden Buck equation, the saturation vapor pressure of water is found to be about . The time constant for drying a red oak board at 150°F is then days, which is the time required to reduce the moisture content to 1/e = 37% of its initial deviation from equilibrium. If the relative humidity is 0.50, then using the Hailwood-Horrobin equation the moisture content of the wood at equilibrium is about 7.4%. The time to reduce the lumber from 85% moisture content to 25% moisture content is then about 4.5 days. Higher temperatures will yield faster drying times, but they will also create greater stresses in the wood due because the moisture gradient will be larger. For firewood, this is not an issue but for woodworking purposes, high stresses will cause the wood to crack and be unusable. Normal drying times to obtain minimal seasoning checks (cracks) in 25mm (1inch or 4/4 lumber) Red Oak ranges from 22 to 30 days, and in 8/4, (50mm or 2inch) it will range from 65 to 90 days.
Methods of drying timber
Broadly, there are two methods by which timber can be dried:
natural drying or air drying
artificial drying
Air drying
Air drying is the drying of timber by exposing it to the air. The technique of air drying consists mainly of making a stack of sawn timber (with the layers of boards separated by stickers) on raised foundations, in a clean, cool, dry and shady place. Rate of drying largely depends on climatic conditions, and on the air movement (exposure to the wind). For successful air drying, a continuous and uniform flow of air throughout the pile of the timber needs to be arranged (Desch and Dinwoodie, 1996).
The rate of loss of moisture can be controlled by coating the planks with any substance that is relatively impermeable to moisture; ordinary mineral oil is usually quite effective. Coating the ends of logs with oil or thick paint improves their quality upon drying. Wrapping planks or logs in materials which will allow some movement of moisture, generally works very well provided the wood is first treated against fungal infection by coating in petrol/gasoline or oil. Mineral oil will generally not soak in more than 1–2 mm below the surface and is easily removed by planing when the timber is suitably dry.
Benefits: It can be less expensive to use this drying method (there are still costs associated with storing the wood, and with the slower process of getting the wood to market), and air drying often produces a higher quality, more easily workable wood than with kiln drying.
Drawbacks: Depending on the climate, it takes several months to a number of years to air-dry the wood.
Kiln drying
The process of artificial or 'oven' drying consists basically of introducing heat. This may be achieved directly, using natural gas and/or electricity, or indirectly, through steam-heated heat exchangers. Solar energy is also an option. In the process, deliberate control of temperature, relative humidity and air circulation creates variable conditions to achieve specific drying profiles. To achieve this, the timber is stacked in chambers that are fitted with equipment to control atmospheric temperature, relative humidity and circulation rate (Walker et al., 1993; Desch and Dinwoodie, 1996).
Chamber drying provides a means of overcoming the limitations imposed by erratic weather conditions. With kiln drying, as is the case with air drying, unsaturated air is used as the drying medium. Almost all commercial timbers of the world are dried in industrial kilns. A comparison of air drying, conventional kiln and solar drying is given below:
Timber can be dried to any desired low moisture content by conventional or solar kiln drying, but in air drying, moisture contents of less than 18% are difficult to attain for most locations.
The drying times are considerably less in conventional kiln drying than in solar kiln drying, followed by air drying.
This means that if capital outlay is involved, this capital sits for a longer time when air drying is used. On the other hand, installing, operating and maintaining an industrial kiln is expensive.
In addition, wood that is being air dried takes up space, which could also cost money.
In air drying, there is little control over the drying conditions, so drying rates cannot be controlled.
The temperatures employed in kiln drying typically kill all the fungi and insects in the wood if a maximum dry-bulb temperature of above 60 °C is used for the drying schedule. This is not guaranteed in air drying.
If air drying is done improperly (exposed to the sun), the rate of drying may be overly rapid in the dry summer months, causing cracking and splitting, and too slow during the cold winter months.
Significant advantages of conventional kiln drying include higher throughput and better control of the final moisture content. Conventional kilns and solar drying both enable wood to be dried to any moisture content regardless of weather conditions. For most large-scale drying operations solar and conventional kiln drying are more efficient than air drying.
Compartment-type kilns are most commonly used in timber companies. A compartment kiln is filled with a static batch of timber through which air is circulated. In these types of kiln, the timber remains stationary. The drying conditions are successively varied according to the type of timber being dried. This drying method is well suited to the needs of timber companies, which have to dry timbers of varied species and thickness, including refractory hardwoods that are more liable than other species to check and split.
The main elements of chamber drying are:
Construction materials The chambers are generally built of brick masonry, or hollow cement-concrete slabs. Sheet metal or prefabricated aluminium in a double-walled construction with sandwiched thermal insulation, such as glass wool or polyurethane foams, are materials that are also used in some modern timber ovens. However, brick masonry chambers, with lime and (mortar) plaster on the inside and painted with impermeable coatings, are used widely and have been found to be satisfactory for many applications.
Heating Heating is usually carried out by steam heat exchangers and pipes of various configurations (e.g. plain, or finned (transverse or longitudinal) tubes) or by large flue pipes through which hot gases from a wood-burning furnace are passed. Only occasionally is electricity or gas employed for heating.
Humidification Humidification is commonly accomplished by introducing live steam into the kiln through a steam spray pipe. In order to limit and control the humidity of the air when large quantities of moisture are being rapidly evaporated from the timber, there is normally a provision for ventilation of the chamber in all types of kilns.
Air circulation Air circulation is the means for carrying the heat to and the moisture away from all parts of a load. Forced circulation kilns are most common, where the air is circulated by means of fans or blowers, which may be installed outside the kiln chamber (external fan kiln) or inside it (internal fan kiln).
Throughout the process, it is necessary to keep close control of the moisture content using a moisture meter system in order to reduce over-drying and allow operators to know when to pull the charge. Preferably, this in-kiln moisture meter will have an auto-shutoff feature.
Kiln drying schedules
Satisfactory kiln drying can usually be accomplished by regulating the temperature and humidity of the circulating air to control the moisture content of the lumber at any given time. This condition is achieved by applying kiln-drying schedules. The desired objective of an appropriate schedule is to ensure drying lumber at the fastest possible rate without causing objectionable degrade. The following factors have a considerable bearing on the schedules.
The species Variations in anatomical, physical, and mechanical properties between species affect drying times and overall results.
The thickness of the lumber Drying time is inversely related to thickness and, to some extent, the width of the lumber.
Whether the lumber boards are quarter-sawn, flat-sawn, or bastard-sawn (mixed-sawn) Sawing pattern influences the distortion due to shrinkage anisotropy.
Permissible drying degrade Aggressive drying schedules can cause timber to crack and distort.
Intended use of timber Mechanical and aesthetic requirements will necessitate different moisture targets depending on the intended use.
Considering each of the factors, no one schedule is necessarily appropriate, even for similar loads of the same species. This is why there is so much timber drying research focused on the development of effective drying schedules.
Dehumidification kiln
A dehumidification chamber can be an unvented system (closed loop) or a partially vented system which uses a heat pump to condense moisture from the air using the cold side of the refrigeration process (evaporator.) The heat thus gathered is sent to the hot side of the refrigeration process (condenser) to re-heat the air and returns this drier and warmer air inside the kiln. Fans blow the air through the piles as in a normal kiln. These kilns traditionally operate from 100 °F to 160 °F and use about half the energy of a conventional kiln.
Vacuum kiln
These kilns can be the fastest to dry and most efficient with energy usage. In a vacuum, water boils at a lower temperature. In addition to increased speed, a vacuum kiln can also produce an improved quality in the wood.
Low ambient pressure does lower the boiling point of water but the amount of energy required to convert the liquid to vapor is the same. Savings come from not being required to heat a huge building and not being required to vent the heat while lowering humidity.
Since all free water can be removed at below 115 °F, quality is improved.
While conventional drying uses warm, dry air to skim water off the surface, vacuum kilns can boil water from within the wood. This enables a good vacuum kiln to dry very thick wood very quickly. It is possible to dry 12/4 Red Oak fresh off the saw to 7% in 11 days.
Since wood is dried with a vapor gradient - vapor pressure to ambient pressure - humidity can be kept very high. Because of this, a good vacuum kiln can dry 4.5" thick White Oak fresh off the saw to 8% in less than a month, a feat that was previously thought to be impossible.
Solar kiln
A solar kiln is a cross between kiln drying and air drying. These kilns are generally a greenhouse with a high-temperature fan and either vents or a condensing system. Solar kilns are slower and variable due to the weather, but are low cost.
Water seasoning
Immersion in running water quickly removes sap and then the wood is air dried. "...it reduces the elasticity and durability of the wood and also makes it brittle." But there are competing perspectives, e.g., "Duhamel, who made many experiments on this important subject, states, that timber for the joiner's use is best put in water for some time, and afterwards dried; as it renders the timber less liable to warp and crack in drying; but, he adds, 'where strength is required it ought not to be put in water.'"
Boiling or steam seasoning
Submersion in boiling water or the application of steam speeds the drying of wood. This method is said to cause less shrinkage "… but it is expensive to use, and reduces the strength and elasticity of the timber."
Chemical or salt seasoning
Salt seasoning is the submersion of wood in a solution of urea, sodium nitrate, all of which act as dehydrating agents. Then the wood is air dried.
Electrical seasoning
Electrical seasoning involves running an electric current through the lumber causing heat to be generated and drying the wood. This method is expensive but is fast and uniform quality.
Freeze drying
Freeze drying is accomplished by lowering the pressure in a chamber containing the wood to a few millibars, while lowering the temperature of the chamber to below the eutectic point of the material. Heat is typically added slowly to the material to allow the water contained in the wood to sublimate directly into vapor, and be deposited on the sides of the vacuum chamber or in the cold trap through which the chamber is evacuated. Freeze drying through sublimation typically takes about 10 times the energy that is taken through evaporation of water by heat. In practice, freeze drying of wood can be accomplished by placing room temperature wood in a vacuum chamber that can be chilled to -30 degrees C or lower, evacuating the chamber to a few millibars, and at the same time cooling the chamber to a freezing temperature. The latent heat of the ice in the wood will come out through the water vapor, which will condense as ice on the inside of the chamber. After a few hours under vacuum and freezing conditions, the chamber is returned to normal pressure, the wood removed and bagged in plastic to keep water from condensing on it, and allowed to return to room temperature over a few hours to a day. The cycle is then repeated, each time the latent heat in the wood is removed through the water content in the wood sublimating and/or evaporating and condensing on the sides of the container and in the cold trap. The cycles are repeated until the moisture content of the wood is at a pre-determined acceptable level. Instead of cycling the wood in the chamber, heat can be added to the wood at a rate that matches the rate of sublimation of ice in the wood to water vapor, which is deposited on the inside of the chamber or in the cold trap. An advantage of freeze drying wood is that the form of the wood is maintained, and shrinkage does not typically occur. Shrinkage will occur over time after the wood is freeze dried, but this typically will not cause defects in the wood.
Drying defects
Drying defects are the most common form of degrade in timber, next to natural problems such as knots (Desch and Dinwoodie, 1996).
There are two types of drying defects, although some defects involve both causes:
Defects from shrinkage anisotropy, resulting in warping: cupping, bowing, twisting, crooking, spring and diamonding.
Defects from uneven drying, resulting in the rupture of the wood tissue, such as checks (surface, end and internal), end splits, honey-combing and case hardening. Collapse, often shown as corrugation, or so-called washboarding of the wood surface, may also occur (Innes, 1996). Collapse is a defect that results from the physical flattening of fibres to above the fibre saturation point and is thus not a form of shrinkage anisotropy.
The standard organizations in Australia and New Zealand (AS/NZS 4787, 2001) have developed a standard for timber quality. The five measures of drying quality include:
moisture content gradient and presence of residual drying stress (case-hardening);
surface, internal and end checks;
collapse;
distortions;
discolouration caused by drying.
Wood-drying kiln
A variety of wood drying kiln technologies exist today: conventional, dehumidification, solar, vacuum and radio frequency.
Conventional wood dry kilns (Rasmussen, 1988) are either package-type (sideloader) or track-type (tram) construction. Most hardwood lumber kilns are sideloader kilns in which fork trucks are used to load lumber packages into the kiln. Most softwood lumber kilns are track types in which lumber packages are loaded on kiln/track cars for loading the kiln.
Modern high-temperature, high-air-velocity conventional kilns can typically dry green lumber in 10 hours down to a moisture content of 18%. However, 1-inch-thick green Red Oak requires about 28 days to dry down to a moisture content of 8%.
Heat is typically introduced via steam running through fin/tube heat exchangers controlled by on/off pneumatic valves. Less common are proportional pneumatic valves or even various electrical actuators. Humidity is removed via a system of vents, the specific layout of which are usually particular to a given manufacturer. In general, cool dry air is introduced at one end of the kiln while warm moist air is expelled at the other. Hardwood conventional kilns also require the introduction of humidity via either steam spray or cold water misting systems to keep the relative humidity inside the kiln from dropping too low during the drying cycle. Fan directions are typically reversed periodically to ensure even drying of larger kiln charges.
Most softwood lumber kilns operate below temperature. Hardwood lumber kiln drying schedules typically keep the dry bulb temperature below . Difficult-to-dry species might not exceed .
Dehumidification kilns are very similar to conventional kilns in basic construction. Drying times are usually comparable. Heat is primarily supplied by an integral dehumidification unit which also serves to remove humidity. Auxiliary heat is often provided early in the schedule where the heat required may exceed the heat generated by the DH unit.
Solar kilns are conventional kilns, typically built by hobbyists to keep initial investment costs low. Heat is provided via solar radiation, while internal air circulation is typically passive.
In 1949 a Chicago company introduced a wood drying kiln that used infrared lamps that they claimed reduced the standard drying time from 14 days to 45 minutes.
Newer wood drying technologies have included the use of reduced atmospheric pressure to attempt to speed up the drying process. A variety of vacuum technologies exist, varying primarily in the method heat is introduced into the wood charge. Hot water platten vacuum kilns use aluminum heating plates with the water circulating within as the heat source, and typically operate at significantly reduced absolute pressure. Discontinuous and SSV (super-heated steam) use atmosphere to introduce heat into the kiln charge. Discontinuous technology allows the entire kiln charge to come up to full atmospheric pressure, the air in the chamber is then heated, and finally vacuum is pulled. SSV run at partial atmospheres (typically around 1/3 of full atmospheric pressure) in a hybrid of vacuum and conventional kiln technology (SSV kilns are significantly more popular in Europe where the locally harvested wood is easier to dry versus species found in North America). RF/V (radio frequency + vacuum) kilns use microwave radiation to heat the kiln charge, and typically have the highest operating cost due to the heat of vaporization being provided by electricity rather than local fossil fuel or waste wood sources.
Valid economic studies of different wood drying technologies are based on the total energy, capital, insurance/risk, environmental impacts, labor, maintenance, and product degrade costs for the task of removing water from the wood fiber. These costs (which can be a significant part of the entire plant costs) involve the differential impact of the presence of drying equipment in a specific plant. An example of this is that every piece of equipment (in a lumber manufacturing plant) from the green trimmer to the infeed system at the planer mill is the "drying system". Since thousands of different types of wood products manufacturing plants exist around the globe, and may be integrated (lumber, plywood, paper, etc.) or stand alone (lumber only), the true costs of the drying system can only be determined when comparing the total plant costs and risks with and without drying.
The total (harmful) air emissions produced by wood kilns, including their heat source, can be significant. Typically, the higher the temperature the kiln operates at, the larger amount of emissions are produced (per pound of water removed). This is especially true in the drying of thin veneers and high-temperature drying of softwoods.
OSHA Standards regarding Dry Kiln Facilities
1910.265(f)(3)(i)(a): Main kiln doors shall be provided with a method of holding them open while kiln is being loaded.
1910.265(f)(3)(i)(b): Counterweights on vertical lift doors shall be boxed or otherwise guarded.
1910.265(f)(3)(i)(c): Adequate means shall be provided to firmly secure main doors, when they are disengaged from carriers and hangers, to prevent toppling.
1910.265(f)(3)(ii)(a): If operating procedures require access to kilns, kilns shall be provided with escape doors that operate easily from the inside, swing in the direction of exit, and are located in or near the main door at the end of the passageway.
1910.265(f)(3)(ii)(b): Escape doors shall be of adequate height and width to accommodate an average size man.
1910.265(f)(4): Pits. Pits shall be well ventilated, drained, and lighted, and shall be large enough to safely accommodate the kiln operator together with operating devices such as valves, dampers, damper rods, and traps.
| Technology | Building materials | null |
5827404 | https://en.wikipedia.org/wiki/Israel%20Space%20Agency | Israel Space Agency | The Israel Space Agency (ISA; , Sokhnut heKhalal haYisraelit) is a governmental body, a part of Israel's Ministry of Science and Technology, that coordinates all Israeli space research programs with scientific and commercial goals.
The agency was founded by the theoretical physicist Yuval Ne'eman in 1983 to replace the National Committee for Space Research, which had established in 1960 to set up the initial infrastructure required for space missions. The agency is currently headed by Chairperson Dan Blumberg and Director General . Today, Israel is the smallest country with indigenous launch capabilities, as well as the smallest to have a space agency.
History
Space research in Israel has history dating to the late 1950s.
NCSR and foreign reliance
The Israeli Space Agency originated as a university-based research project from Tel Aviv University in the early 1960s. In 1960, the National Committee for Space Research (NCSR) was formed by the Israel Academy of Sciences and Humanities. The committee was formed to increase research activities across the academic communities in Israel. While at the time establishing a space program was not particularly one of its goals, during the 1960s through the late 1970s, the committee developed the infrastructure needed for research and development in space exploration and sciences. One of the NCSR's earliest achievements took place in 1961 with the launch of its first two-stage rocket.
Following political tension with Egypt and Syria, reconnaissance flights became evermore difficult. In 1979, a satellite program was proposed followed by a year of feasibility study. The study was completed by late 1980; Saguy requested from prime minister Menachem Begin that the project proceed to its next phase. In 1982, a new recommendation called "Ofeq Program" was submitted for developing an observation satellite. The program included timelines, planning for a ground station, budget estimates, and personnel requirements. The primary goal was to develop a satellite program without relying on any foreign know-how, to allow flexibility and creativity. The launcher would be developed by Malam, the two engines by the Israel Military Industries (IMI), and the third by Rafael.
Establishment of a space agency
At the end of 1982 it was decided during a closed-door meeting to establish an Israeli space agency. The decision was made by PM Menachem Begin, Defense minister Ariel Sharon, and former director Aharon Beit Halahmi. The initial goal was to pursue the program to develop to Ofeq and the Shavit launchers. In January 1983, the Israeli government authorized the Minister of Science and Technology Yuval Ne'eman to establish an Israeli Space Agency with the goal of advancing Israel's space program, unlike the NCSR which was primarily used for feasibility and infrastructure studies. In July 1983, the ISA was officially founded in Tel Aviv to coordinate the nation's space program in affiliation to the Ministry of Science, Culture & Sport, and Dror Sadeh was nominated to be its Director-General.
In 1982, the Israel Space Agency was created, which is responsible for coordinating the space program of Israel. The following year, the new head of military intelligence, Ehud Barak, suspended all work, and advocated the winding down of all projects, and planned to transfer the released financial resources to fulfill more priority tasks, he said. But his opinion did not become dominant in the defense department, and in 1984, Defense Minister Moshe Arens insisted on resuming the program.
In 1984 the National Space Knowledge Center was established in cooperation with Israel Aircraft Industries; a contract was signed between IAI and Ministry of Defence for the development of the needed infrastructure and of Israel's first observation satellite. This came to fruition in 1988 when Israel launched the first in a series of Ofeq satellites and thus became one of only a few nations in the world possessing an indigenous space launching capability (see timeline of first orbital launches by country). The project management at Israel Aircraft Industries was headed for many years by Moshe Bar-Lev.
Vision
The agency vision as defined by the guiding committee on July 27, 2005, states:
"Space research and exploration is an essential instrument for the defense of life on Earth; the lever for technological progress; the key to existing in a modern society; essential for developing an economy based on knowledge; and the central attraction for scientific and qualified human resources."
The vision is "to preserve and broaden the comparative advantage of Israel and to place it among the group of leading countries in the space research and exploration area."
The main goals for vision realization are:
To build and to support satellite systems for space research and for Earth research from space.
To develop technologies, knowledge and scientific infrastructure (including laboratories and human resources) required for space research.
To promote international cooperation in space research and exploration, and for strengthening the national interests of Israel.
To promote ties between Israeli society, space research, and exploration.
ISA has signed cooperation agreements with the space agencies of: United States (NASA), France (CNES), Canada (CSA), India (ISRO), Italy (ASI), Germany (DLR), Ukraine (NSAU), Russia (RKA), Netherlands (NIVR) and Brazil (AEB).
Budget
In 2010, the budget of the Israel Space Agency was increased to US $80 million to boost the agency's space activities in research and development. The budget does not include launch vehicle development and most satellite programs. Typically such programs get funded on a project-by-project basis. For example, Project Venus, a cooperative program by Israel (ISA) and France (CNES) which is set to launch in 2014 has a $50 million budget. The Spaceborne Hyperspectral Applicative Land and Ocean Mission (SHALOM), a joint mission by Israel (ISA) and the Italy Space Agency (ASI) has a budget of $116 million.
The budget allocated annually for the Israeli military program as well as commercial programs are managed on different budgets.
Satellite programs
The Israel Space Agency has had a long history of satellite programs both for reconnaissance and commercial purposes. Its first satellite, the Ofeq-1 was launched on September 19, 1988, from Palmachim Airbase in Israel. Since the launching of that first satellite, Israel has developed into a significant player in the commercial space arena. Today, the ISA satellite launches include:
Ofeq – Series of reconnaissance satellites. The first of these was launched from the Palmachim site on September 19, 1988.
AMOS – Series of communications satellites
Eros – Series of observation satellites
Techsat – Researching satellite launched by the Technion
TechSAR – a SAR-based observation satellite.
Ofeq satellite series
After the successful launch of the Ofeq-1 in 1988, additional satellites were developed. In 1989, the ISA launched the Ofeq-2; in April 1995, it took a leap forward with the launch of Ofeq-3, which carried an advanced electro-optical payload built by Israeli industry for local purposes. Ofeq-3 has been functioning without a hitch. Following a setback with Ofeq-4, Ofeq-5 was successfully launched in May 2002.
To date, twelve such satellites in the Ofeq reconnaissance satellites series were developed and launched to Low Earth Orbit. The most recent, Ofek 16, was launched July 6, 2020.
AMOS satellite series
The AMOS is a series of communications satellites. The AMOS satellites are by the Israel Aerospace Industries and are operated by Spacecom once in orbit. The AMOS-1, the first satellite in the series, was launched on May 16, 1996, using a French-built vehicle. Since then, 4 more satellites were launched. The most recent AMOS satellite is the AMOS-5 which was launched on December 11, 2011.
The AMOS is distinguished for its light weight and sophisticated technology. The AMOS-6 was due to be launched in 2016 in order to replace the AMOS-2 which was then ceasing operation, but it was destroyed in a launchpad explosion.
EROS satellite series
The Earth Resources Observation Satellite (EROS) is a series of commercial observation satellites. The first satellite, the EROS A, was launched on December 5, 2000, from Svobodny Launch Complex using a Russian Start-1. A second satellite, the EROS B, was launched on April 25, 2006.
The EROS series are set to be launched once every 6 to 8 years. The EROS C is set to be launched in late 2014.
TechSAR Satellite
The TechSAR satellite is a reconnaissance satellite equipped with a synthetic aperture radar. The satellite is designed penetrate thick clouds by being fitted with a large dish-like antenna to transmit and receive radar signals. The satellite was successfully launched on January 21, 2008.
Launch capabilities
The Israel Space Agency is one of only seven countries that both build their own satellites and launch their own launchers. The Shavit is a space launch vehicle capable of sending payload into low Earth orbit. The Shavit launcher has been used to send every Ofeq satellite to date.
The development of the Shavit began in 1983 and its operational capabilities were proven on three successful launches of the Ofek satellites on September 19, 1988; April 3, 1990; and April 5, 1995. The Shavit launchers allows low-cost and high-reliability launch of micro/mini satellites to a low Earth orbit. The Shavit launcher is developed by Malam factory, one of four factories in the IAI Electronics Group. The factory is very experienced in development, assembling, testing and operating system for use in space.
The Shavit is a triple-stage launcher solid propellant booster based on the 2-stage Jericho-II ballistic missile. The first and second stage engines are manufactured by Ta'as, and use solid fuel. The third stage engines are manufactured by Rafael Advanced Defense Systems. The next generation Shavit rockets, now called the Shavit-2 are being developed. The Shavit-2 is said to be made available for commercial launches in the near future.
Palmachim Spaceport
The Israel Space Agency and the Israeli Air Force (IAF) jointly operate the spaceport located south of Palmachim Airbase. The IAF with its 151 Squadron for missile testing uses it to test its Jericho and Arrow missiles.
Due to Israel's geographic location and hostile relations with surrounding countries, launches take off due west, over the Mediterranean Sea. This is done in order to avoid flying over hostile territories. This is also to prevent possible debris from falling above populated areas. This limitation imposes a penalty of roughly 30% on its lifting capabilities.
Some of the recent ISA launches include:
June 11, 2007 – Ofeq 7 satellite
June 22, 2010 – Ofeq 9 satellite
April 10, 2014 – Ofeq 10 satellite
September 13, 2016 – Ofeq 11 satellite
July 6, 2020 – Ofeq 16 satellite
Space research
The ISA is responsible for funding a large set of university research projects and governmental projects. The goal is to boost space related research and development in the academic arena. The Israel Space Agency strives to promote space research and space technology development as a part of the effort to promote the Israeli scientific research.
National Knowledge Center on NEOs
The ISA and the Ministry of Science and Technology formed and operated the National Knowledge Center on Near Earth Objects at Tel Aviv University in order to study minor bodies in the Solar System. The goal is to map the objects which pose a threat to Earth and to find a solution to eliminate them.
This center is headed by Dr. Noah Brosch from Tel Aviv University and operates the three telescopes of the Wise Observatory. For the purpose of NEO studies, a special wide-field 0.46-m telescope was acquired by Tel Aviv University and is operating at the Wise Observatory located near Mitzpe Ramon. The telescope facilitated the discovery of several tens of new asteroids and is now intensely used to study asteroid rotation and asteroid pair properties. Following the cessation of funding from ISA, the asteroid studies are continuing at the Wise Observatory with no specific financial support.
Israel Cosmic Ray Center
The Israel Cosmic Ray Center (ICRC) was established in November 1997 with support from the Israel Space Agency and with affiliation to Tel Aviv University and the Technion. The center is located on Mount Hermon and headed by Gideon Bela and Lev Dorman from the Tel Aviv University.
The goal of the center is to monitor and forecast dangerous meteorological and space phenomena. This includes solar radiation storms and shockwaves between very powerful stars creating magnetic storms. These phenomena can endanger electronic systems in satellites and space shuttles, the astronauts' health, electronic and navigational systems in aircraft flying in extremely high altitudes and ground power systems.
The Emilio Segre' Observatory (ESOI) is a cosmic radiation observatory prepared via a joint collaboration with Italy. It was transported to Mount Hermon, Israel in June 1998. The observatory was named after the Italian physicist and Nobel laureate Emilio G. Segrè.
EOSDIS ISA-MEIDA
The Israel Space Agency - Middle East Interactive Data Archive (ISA-MEIDA) is an Israeli node for NASA's EOSDIS (Earth Observing System Data and Information System). The Node was established in October 1996 as a part of the cooperation agreement between the director of NASA and the director of ISA. It is the only team in the country which focuses on collecting and preserving environmental information in Israel. The ISA-MEIDA was established in order to create and maintain an Earth observing data center available through the Internet to the research community and to the general public free of charge.
This research, headed by Pinhas Alpert, is funded by the ISA as a part of the Space Scientific and Technological Infrastructure Development program. The node is integrated with NASA's Global Information System which includes Earth Observing Science Data and Information System. It also contains Remote Sensing Data from NASA and NOAA satellites, and data from other sources such as meteorological RADAR and readings from meteorological stations. Information collected by the system is critical for environmental and water related research, especially in light of the drastic global environmental changes and global warming.
Tel Aviv University Ultraviolet Explorer
The Tel Aviv University Ultraviolet Explorer (TAUVEX) is a space telescope array conceived by Noah Brosch of Tel Aviv University. It was eventually designed and constructed for the university by El-Op. TAUVEX is a cluster of three bore-sighted 20-cm telescopes for observation in the ultra-violet funded by the Israel Space Agency. It was due to be launched by India on board the ISRO satellite GSAT-4 satellite but due to a mis-match between the capabilities of the launched and the mass of the satellite, ISRO decided unilaterally to remove TAUVEX from GSAT-4. The telescope ended up in a limbo of constant delays. In 2012 ISA decided to terminate the TAUVEX project.
TechSat-Gurwin Microsatellite
The Gurwin TechSat was one of the world's first Microsatellites to be designed, built and launched by students. The satellite, developed by the Technion, was launched on July 10, 1998. The satellite was a great success, operating for over 11 years - setting a world record for the longest university satellite mission. The mission was completed in April 2010.
Sloshsat-FLEVO
The Sloshsat-FLEVO was a satellite launched for the study of liquid sloshing phenomena in space. The satellite was developed in collaboration between the Dutch NLR labs, Fokker Space, and Rafael. The satellite is helped by the Israel Space Agency and the European Space Agency (ESA). The Israel Space Agency initiated the project and partly funded it by supplying the sub-propulsion system.
The propulsion system, funded by the ISA, is fueled by cold gas and is constructed of four high-pressure canisters with pyrotechnic activation. The system can produce both linear acceleration and torque needed to perform slosh research. It is specifically designed to overcome the severe volume constraints encountered in typical satellites.
Ongoing development
The Israel Space Agency is currently involved with multiple satellites, space telescopes, and microsatellites.
VENμS
The Vegetation and Environment monitoring on a New Micro-Satellite (VENμS) is a satellite to be used for Earth observation using a superspectral sensor, dedicated to vegetation monitoring. It is the first cooperation between Israel (ISA) and France (CNES). VENμS scientific objective is "the provision of data for scientific studies dealing with the monitoring, analysis, and modeling of land surface functioning under the influences of environmental factors as well as human activities."
The Israel Space Agency is responsible for the spacecraft, the launcher interface, and for the satellite control center. The Centre national d'études spatiales is responsible for supplying the superspectral camera and the science mission center. The satellite was launched in August 2017.
At World Space Week in Dubai, in October 2021, Israel and the UAE ministers of science and technology announced plans for cooperation on the mission. They plan to examine phenomena related to Earth resources, precision agriculture, desertification, and monitoring of bodies of water, and climate change– issues common both to Israel and the Emirates.
OPsat
OPsat is a next generation high resolution optical observation satellite for reconnaissance purposes. It is designed to be a satellite capable of detecting objects of about 50 centimeters in diameter. It will be equipped with a camera with CCD/TDI sensors, producing both panchromatic imagery at a very high resolution and multispectral imagery at a medium resolution. The satellite is set to orbit in a Sun-synchronous orbit. It is expected to have a lifespan of roughly 10 years.
ULTRASAT
The ULTRASAT satellite is a wide field (~1000 deg2) transient explorer space telescope mission. ULTRASAT is planned to have eight telescopes equipped with CCD cameras and reflective filters. It is set to have a sensitivity 10 times lower than GALEX but a field of view more than 1000 times larger. It is also planned to have a detection rate for transient in the UV of more than 30 times greater than that of GALEX. The ULTRASAT is planned to be developed in just 3–4 years and relatively low cost. ULTRASAT will observe a large patch of sky, more than 200 square degrees, alternating every six months between the southern and northern hemisphere. The satellite will orbit the Earth from an altitude of about 300 km above the geosynchronous orbit, getting a ‘ride’ as a secondary payload in the fairing of the rocket carrying a communications satellite. A joint American-Israeli proposal for this project was submitted to NASA by a team from Caltech/JPL (Jet Propulsion Laboratory), the Weizmann Institute of Science and Israel Aerospace Industries (IAI). The Israeli contribution will be funded by the Israel Space Agency and the launch is scheduled for early 2026.
INSAT-1 and INSAT-2
The INSAT-1 and the INSAT-2 are two nano-satellites which were planned and developed by the Israeli Nano Satellite Association. Their purpose was to serve as a technology demonstrators.
The INSAT-1 was set to carry a miniature atomic clock and a GPS receiver. By 2021, the mission was delayed and they were never launched.
SHALOM Project
The SHALOM Project (Spaceborne Hyperspectral Applicative Land and Ocean Mission) is a joint project with the Italian Space Agency (ASI) announced on November 23, 2010. The project involves the design and development of two hyperspectral Earth observation satellites. They are set to follow the same orbit as the Cosmo SkyMed. It will integrate radar observations with observations in the visible infrared and ultraviolet. In October 2015 a memorandum of understanding was signed, and the system is slated to become fully operational in 2021. The project is expected to cost over $200 million, with the cost being split evenly between the two countries.
SAMSON
Space Autonomous Mission of Swarming & Geolocating Nano-Satellites (SAMSON) is a project initiated by the Technion's Asher Space Research Institute. SAMSON consists of three nano-satellites in formation flying to demonstrate high precision geo-location of civilian signals from the ground for rescue purposes. SAMSON is planned to implement and demonstrate Technion-developed formation-flying algorithms using the nano-satellite's propulsion system. This is a student project with technical help from multiple partners in the industry such as RAFAEL.
Matroshka AstroRad Radiation Experiment - Artemis 1
The Matroshka AstroRad Radiation Experiment (MARE) is a collaboration between the ISA, NASA, and the German Aerospace Center (DLR), in cooperation with Lockheed Martin and StemRad. MARE will take place aboard the Artemis 1 Orion Spacecraft, which is expected to launch from Kennedy Space Center between 2020 and 2021. The radiation environment beyond Earth's protective atmosphere is a significant impediment to human space travel. The AstroRad is a protective vest that provides mobile shielding from high energy radiation. MARE offers a unique opportunity to measure tissue dose deposition and test the effectiveness of the AstroRad vest when exposed to the harsh radiation beyond low Earth orbit. The crew compartment of the uncrewed Artemis 1 Orion spacecraft will include two female anthropomorphic phantoms that will be exposed to the intense radiation environment during the spacecraft's lunar orbit. One phantom will be shielded by the AstroRad vest and the other will be left unprotected as a control. The phantoms, designed by DLR, precisely measure radiation exposure not only at the surface of the body but also at the locations corresponding to sensitive internal organs and tissues in humans.
Crewed programs
While the ISA has not yet made a human spaceflight of its own, since the mid-1990s there has been a cooperation agreement between the Israel Space Agency and NASA which has resulted in one Israeli astronaut on a NASA mission to date.
Ilan Ramon
Ilan Ramon was Israel's first astronaut. Ramon was the Space Shuttle payload specialist on board the fatal STS-107 mission of Space Shuttle Columbia, in which he and the six other crew members were killed in a re-entry disaster over southern Texas. Ramon had been selected as a Payload Specialist in 1997 and trained at the Johnson Space Center, Houston, Texas, from 1998 until 2003, for a mission with a payload that included a multispectral camera for recording desert aerosol (dust).
Experiments via NASA Space Shuttles
The ISA has conducted a number of experiments, both crewed and uncrewed, in collaboration with NASA using the Space Shuttle. In October 1996, NASA and ISA signed an agreement for joint cooperation in the peaceful use of space - an agreement designed to develop cooperative programs of mutual interest between the two nations.
Israeli Space Agency Investigation About Hornets
The Israeli Space Agency Investigation About Hornets (ISAIAH) was a project from Tel Aviv University that was initiated in the early 1990s to explorer the effects of near-zero gravity on oriental hornets, their physical and physiological development and their nest-building instincts. The flight hardware and measuring instruments were commissioned by the ISA and built by ISI. The hope of the mission was to discover ways to prevent astronauts from suffering headaches, nausea, and vomiting during the missions.
In 1992, 230 Oriental hornets, a flight hardware and measuring instruments were packed onto the Space Shuttle Endeavour mission STS-47. The Oriental Hornets used in the experiment were capable of building combs in the direction of the gravitational vector and detecting gravitational force changes in real time.
Early development of embryos in microgravity
The early development of mice embryos in microgravity was a Hebrew University experiment designed to determine if mice embryo cells could develop normally in microgravity conditions in space. The outcome of the experiment was expected to help the understanding of early embryo cell development which would provide an insight into the possibilities of human reproduction in space. The payload was launched in 1996 on board Space Shuttle Columbia mission STS-80.
On the same mission, a second experiment investigating the growth of osteoblast cells in microgravity environment was also conducted. During STS-80 osteoblast cultures were grown in microgravity, using specialized hardware. The space-grown osteoblastic cells were then compared with the Earth-grown osteoblastic cells. The research revealed numerous changes between the two cells. The microgravity cells showed lower proliferation rate, a lower metabolism and an altered cell structure.
A continuation from the osteoblast cells experiment was later expanded upon by astronaut John Glenn on board Space Shuttle Discovery, STS-95. The tests included the thinning effect of space on mouse bones. Additionally the process of calcium loss in the mouse bone was later compared to what happened in Glenn's body and effectiveness of a calcium-vitamin D supplement against osteoporosis.
MEIDEX
The 2003 Mediterranean Israeli Dust Experiment (MEIDEX) was a Dust Experiment initiated by the ISA. The experiment was planned by a team from the Department of Geophysics and Planetary Sciences in Tel Aviv University. The objective was to study the temporal and spatial distribution and physical properties of atmospheric desert dust over North Africa, the Mediterranean and the Atlantic Saharan regions. The aim was achieved by a remote sensing experiment operated by the astronaut Ilan Ramon aboard the Space Shuttle Columbia.
Moon mission
The Israeli Space Agency has sponsored the first Israeli effort to land a spacecraft on the Moon. Additionally the ISA has been collaborating with NASA about future lunar research programs.
Israel Network for Lunar Science and Exploration
The Israel Network for Lunar Science and Exploration (INLSE) program was established by the Israeli Space Agency part of an international effort to study the Moon and the Solar System. In January 2010 a joint declaration by NASA and Israel Space Agency was signed making Israel a member of the NASA Center for Moon Research and promote cooperation between the two agencies. The INLSE hopes to bring its technical and engineering expertise for the sake of advancing the broad goals of lunar science at the institute. The agreement offers NASA important research involving lasers, the development of advanced sensors for Solar System research tasks and automatic vehicle navigation.
SpaceIL
The SpaceIL is a non-governmental organization made up of multidisciplinary team of Israeli scientists and space aficionados. The organization was formed to compete in the international Google Lunar X Prize competition.
The SpaceIL team developed a robotic spacecraft, built as a microsatellite, weighing around 500 kg. It was designed to be launched and then land on the Moon bearing the Israeli flag. To win, they must have been the first to launch, fly, and land the spacecraft on the Moon then transmit live video feed back to Earth. The team intended to donate the prize, about $30 million, toward space education, if they won it.
On 11 April 2019, Beresheet crash-landed on the lunar surface. Prior to impact, the probe had been able to take two last photographs: a view of itself against the Moon, and a closer shot of the Moon's surface.
On 13 April 2019, Morris Kahn announced that a new mission, named Beresheet 2 would attempt a second time to land on the Moon. The planned date is 2024. At World Space Week in Dubai, in October 2021, Israel and the UAE ministers of science and technology announced plans for cooperation on the mission. Said Kahn to the Global Investment Forum in Dubai, "It would be wonderful if we could develop a space program that would be a combination of Israel and the Arab world, “I would welcome it – if it fits in with the program the Emirates have. They have an ambitious program.”
Commercial and industry involvement
The Israel Space Agency, Israeli industry, and the academy are all heavily involved in all the different stages of planning, development, construction, launching, and operating of space programs.
The main contractors of the Ofeq and Eros space programs is the IAI Mabat factory in Yehud. The Mabat facility is responsible for the experiment and integration center, the ground monitoring, and the control stations and the remote satellites receiving stations. Many of the high-tech companies are involved in the various space programs in Israel, and in manufacturing sub-systems and components.
The TAUVEX was spun off as DAVID, a small telescope with a resolution of about five meters. It is currently being developed jointly by an Israeli hi-tech firm and a German firm. Many smaller university-related commercial products are being developed across the state.
Industrial groups
In addition to university research, a number of large industrial groups are heavily involved with the Israel Space Agency.
IAI
The Israel Aerospace Industries (IAI) is Israel's prime aerospace and aviation manufacturer, producing aerial systems for both military and civilian usage. IAI entered the space race in the 1990s and has since been responsible for the development of most Israel's civilian and military satellites, particularly the AMOS and Ofeq.
Elbit Systems
El-Op, which merged with Elbit Systems in 2000, is the country's largest research and development company for space-qualified cameras and advanced telescopes that deal with various panchromatic, dual band, multi and super spectral wavelength applications. It is Israel's Center of Excellence for space electro-optics.
Rafael Advanced Defense Systems
Rafael Advanced Defense Systems is the Israeli authority for development of weapons and military technology. They are responsible for most of the ISA's spacecraft propulsion systems.
Agency managers/directors
Dror Sadeh
David Abir
Akiva Bar-Nun 1989 - 1993
Marcel Klein
Aby Har Even (Avi Har-Even) January 1995 - September 2004
Zvi Kaplan September 2004 - September 2011
Menachem Kidron January 2012 - 2015
Avi Blasberger May 2016- August 2021
Uri Oron September 2021
Agency chairperson
Yuval Ne'eman (1983-2005)
Yitzhak Ben Israel (2005-2022)
Dan Blumberg (2022-)
| Technology | Programs and launch sites | null |
7616391 | https://en.wikipedia.org/wiki/Polytrauma | Polytrauma | Polytrauma and multiple trauma are medical terms describing the condition of a person who has been subjected to multiple traumatic injuries, such as a serious head injury in addition to a serious burn. The term is defined via an Injury Severity Score (ISS) equal to or greater than 16. It has become a commonly applied term by US military physicians in describing the seriously injured soldiers returning from Operation Iraqi Freedom in Iraq and Operation Enduring Freedom in Afghanistan. The term is generic, however, and has been in use for a long time for any case involving multiple trauma.
Civilian medicine
In civilian life, polytraumas often are associated with motor vehicle crashes. This is because car crashes often occur at high velocities, causing multiple injuries. On admission to hospital any trauma patient should immediately undergo x-ray diagnosis of their cervical spine, chest, and pelvis, commonly known as a 'trauma series', to ascertain possible life-threatening injuries. (Where available, a CT trauma series for head, neck, thorax, abdomen and pelvis may be the imaging modality of first choice). Examples would be a fractured cervical vertebra, a severely fractured pelvis, or a haemothorax. Once this initial survey is complete, x-rays may be taken of the limbs to assess the possibility of other fractures. It also is quite common in severe trauma for patients to be sent directly to CT or a surgery theatre, if they require emergency treatment.
Extracorporeal membrane oxygenation (ECMO) may be effective in treating some polytrauma patients with pulmonary or cardiopulmonary failure.
Military medicine
Polytrauma often results from blast injuries sustained from improvised explosive devices, or by a hit with a rocket-propelled grenade, with "Improvised explosive devices, blasts, landmines, and fragments account[ing] for 65 percent of combat injuries ...". The combination of high-pressure waves, explosive fragments, and falling debris may produce multiple injuries including brain injury, loss of limbs, burns, fractures, blindness, and hearing loss, with 60 percent of those injured in this way, having some degree of traumatic brain injury.
In some respects, the high incidence of polytrauma in military medicine is, in fact, a sign of medical advancement. In previous wars most soldiers with such multiple injuries simply did not survive, even if quickly transferred into hospital care. Today many polytrauma victims never fully regain their previous physical capacity, and are more susceptible to psychological complications, such as PTSD.
U.S. treatment
As of 2013, there were five rehabilitation centers in the U.S. specialising in polytrauma. They are managed by the United States Department of Veterans Affairs and are located in Minneapolis, Minnesota; Palo Alto, California; Richmond, Virginia; San Antonio, Texas, and Tampa, Florida. In addition to the intensive care, insofar as still required, these hospitals mainly specialize in rehabilitative treatment. In addition the Department of Veterans Affairs has 22 polytrauma network sites, located throughout the country.
Veterans Health Administration (VHA) developed a screening and evaluation process to ensure that OEF/OIF/OND Veterans with TBI are identified, and that they receive appropriate treatments and services. This includes mandatory screening for deployment-related TBI of all OEF/OIF/OND Veterans upon their initial entry into VHA for services. Veterans with positive screens are referred for a comprehensive evaluation by TBI specialists for diagnostic and treatment recommendations.
Based on extensive research, the VA-TBI Screening Tool has revealed high sensitivity and moderate specificity allowing VA to identify symptomatic Veterans and develop an appropriate plan of care. From 2007 to 2015, over 900,000 Veterans have been screened for possible OEF/OIF/OND deployment related TBI. Of those, approximately 20 percent had positive screens and were referred for further evaluation.
Epidemiology
OEF/OIF/OND veterans have a high polytrauma rate. Respectfully, a study exhibited findings with a population of 16,590 OEF/OIF/OND veterans, in which 27.66% met the criteria for poly trauma. Those within this subpopulation were most likely male (92.9%) and White (71.0%). Similar findings in a sample of 2,441,698 OEF/OIF/OND active duty found that the rate of poly trauma was 5.99 per 1,000 individuals. Of those with polytrauma, 52.15% were most likely between the ages of 20–29 years, male (89.93%), White (69.07%), married (64.18%), and enlisted in the Army (74.71%). Furthermore, the rate of polytrauma among a sample of 613,391 OEF/OIF/OND veterans was 6% (36,800). Additional research has concluded that in a selection of 340 OEF/OIF/OND veterans, 42.1% exhibited symptoms of poly trauma.
As of April 2007, the Department of Veterans Affairs has treated more than 350 service members in their inpatient centers.
The treatment and rehabilitative care for polytrauma patients is a very extensive and time-consuming activity. The recommended staffing numbers (FTE = Full Time Equivalent) for six rehabilitation treatment beds are:
0.5 FTE – Physician Discipline FTE Rehabilitation
5.5 FTE – Registered Nurse (1.0 must be CRRN)
4.0 FTE – Licensed Practical Nurse and/or Certified Nursing Assistant
0.5 FTE – Nurse Manager
0.5 FTE – Clinical Case Manager, Admission and Follow-up
1.0 FTE – Social Worker Case Manager
0.5 FTE – Social Worker
1.0 FTE – Speech-Language Pathologist
1.0 FTE – Physical Therapist
1.0 FTE – Occupational Therapist
0.5 FTE – Recreation Therapist
0.5 FTE – Counseling Psychologist
0.5 FTE – Neuropsychologist
In other words, 2.8 people are required full-time (24h), for every patient, often for months, while some care may be required for life.
| Biology and health sciences | Injury: General | Health |
7620308 | https://en.wikipedia.org/wiki/Elementary%20reaction | Elementary reaction | An elementary reaction is a chemical reaction in which one or more chemical species react directly to form products in a single reaction step and with a single transition state. In practice, a reaction is assumed to be elementary if no reaction intermediates have been detected or need to be postulated to describe the reaction on a molecular scale. An apparently elementary reaction may be in fact a stepwise reaction, i.e. a complicated sequence of chemical reactions, with reaction intermediates of variable lifetimes.
In a unimolecular elementary reaction, a molecule dissociates or isomerises to form the products(s)
At constant temperature, the rate of such a reaction is proportional to the concentration of the species
In a bimolecular elementary reaction, two atoms, molecules, ions or radicals, and , react together to form the product(s)
The rate of such a reaction, at constant temperature, is proportional to the product of the concentrations of the species and
The rate expression for an elementary bimolecular reaction is sometimes referred to as the law of mass action as it was first proposed by Guldberg and Waage in 1864. An example of this type of reaction is a cycloaddition reaction.
This rate expression can be derived from first principles by using collision theory for ideal gases. For the case of dilute fluids equivalent results have been obtained from simple probabilistic arguments.
According to collision theory the probability of three chemical species reacting simultaneously with each other in a termolecular elementary reaction is negligible. Hence such termolecular reactions are commonly referred as non-elementary reactions and can be broken down into a more fundamental set of bimolecular reactions, in agreement with the law of mass action. It is not always possible to derive overall reaction schemes, but solutions based on rate equations are often possible in terms of steady-state or Michaelis-Menten approximations.
| Physical sciences | Basics_3 | Chemistry |
7622892 | https://en.wikipedia.org/wiki/Point%20source | Point source | A point source is a single identifiable localized source of something. A point source has a negligible extent, distinguishing it from other source geometries. Sources are called point sources because, in mathematical modeling, these sources can usually be approximated as a mathematical point to simplify analysis.
The actual source need not be physically small if its size is negligible relative to other length scales in the problem. For example, in astronomy, stars are routinely treated as point sources, even though they are in actuality much larger than the Earth.
In three dimensions, the density of something leaving a point source decreases in proportion to the inverse square of the distance from the source, if the distribution is isotropic, and there is no absorption or other loss.
Mathematics
In mathematics, a point source is a singularity from which flux or flow is emanating. Although singularities such as this do not exist in the observable universe, mathematical point sources are often used as approximations to reality in physics and other fields.
Visible electromagnetic radiation (light)
Generally, a source of light can be considered a point source if the resolution of the imaging instrument is too low to resolve the source's apparent size. There are two types and sources of light: a point source and an extended source.
Mathematically an object may be considered a point source if its angular size, , is much smaller than the resolving power of the telescope:
, where is the wavelength of light and is the telescope diameter.
Examples:
Light from a distant star seen through a small telescope
Light passing through a pinhole or other small aperture, viewed from a distance much greater than the size of the hole
Light from a street light in a large-scale study of light pollution or street illumination
Other electromagnetic radiation
Radio wave sources that are smaller than one radio wavelength are also generally treated as point sources. Radio emissions generated by a fixed electrical circuit are usually polarized, producing anisotropic radiation. If the propagating medium is lossless, however, the radiant power in the radio waves at a given distance will still vary as the inverse square of the distance if the angle remains constant to the source polarization.
Gamma ray and X-ray sources may be treated as a point source if sufficiently small. Radiological contamination and nuclear sources are often point sources. This has significance in health physics and radiation protection.
Examples:
Radio antennas are often smaller than one wavelength, even though they are many meters across
Pulsars are treated as point sources when observed using radio telescopes
In nuclear physics, a "hot spot" is a point source of radiation
Sound
Sound is an oscillating pressure wave. As the pressure oscillates up and down, an audio point source acts in turn as a fluid point source and then a fluid point sink. (Such an object does not exist physically, but is often a good simplified model for calculations.)
Examples:
Seismic vibration from a localised seismic experiment searching for oil
Noise pollution from a jet engine in a large-scale study of noise pollution
A loudspeaker may be considered as a point source in a study of the acoustics of airport announcements
A coaxial loudspeaker is designed to work as a point source to allow a wider field for listening.
Ionizing radiation
Point sources are used as a means of calibrating ionizing radiation instruments. They are usually sealed capsules and are most commonly used for gamma, x-ray and beta-measuring instruments.
Heat
In a vacuum, heat escapes as radiation isotropically. If the source remains stationary in a compressible fluid such as air, flow patterns can form around the source due to convection, leading to an anisotropic pattern of heat loss. The most common form of anisotropy is the formation of a thermal plume above the heat source.
Examples:
Geological hotspots on the surface of the Earth which lie at the tops of thermal plumes rising from deep inside the Earth
Plumes of heat studied in thermal pollution tracking.
Fluid
Fluid point sources are commonly used in fluid dynamics and aerodynamics. A point source of fluid is the inverse of a fluid point sink (a point where fluid is removed). Whereas fluid sinks exhibit complex rapidly changing behavior such as is seen in vortices (for example water running into a plug-hole or tornadoes generated at points where air is rising), fluid sources generally produce simple flow patterns, with stationary isotropic point sources generating an expanding sphere of new fluid. If the fluid is moving (such as wind in air or currents in water) a plume is generated from the point source.
Examples:
Air pollution from a power plant flue gas stack in a large-scale analysis of air pollution
Water pollution from an oil refinery wastewater discharge outlet in a large-scale analysis of water pollution
Gas escaping from a pressurized pipe in a laboratory
Smoke is often released from point sources in a wind tunnel in order to create a plume of smoke which highlights the flow of the wind over an object
Smoke from a localized chemical fire can be blown in the wind to form a plume of pollution
Pollution
Sources of various types of pollution are often considered as point sources in large-scale studies of pollution.
| Physical sciences | Electromagnetic radiation | Physics |
662776 | https://en.wikipedia.org/wiki/Turbellaria | Turbellaria | The Turbellaria are one of the traditional sub-divisions of the phylum Platyhelminthes (flatworms), and include all the sub-groups that are not exclusively parasitic. There are about 4,500 species, which range from to large freshwater forms more than long or terrestrial species like Bipalium kewense which can reach in length. All the larger forms are flat with ribbon-like or leaf-like shapes, since their lack of respiratory and circulatory systems means that they have to rely on diffusion for internal transport of metabolites. However, many of the smaller forms are round in cross section. Most are predators, and all live in water or in moist terrestrial environments. Most forms reproduce sexually and with few exceptions all are simultaneous hermaphrodites.
The Acoelomorpha and the genus Xenoturbella were formerly included in the Turbellaria, but are no longer regarded as Platyhelminthes. All the exclusively parasitic Platyhelminthes form a monophyletic group Neodermata, and it is agreed that these are descended from one small sub-group within the free-living Platyhelminthes. Hence the "Turbellaria" as traditionally defined are paraphyletic.
Description
Traditional classifications divide the Platyhelminthes into four groups: Turbellaria and the wholly parasitic Trematoda, Monogenea and Cestoda. In this classification the Turbellaria include the Acoelomorpha (Acoela and Nemertodermatida). The name "Turbellaria" refers to the "whirlpools" of microscopic particles created close to the skins of aquatic species by the movement of their cilia.
Features common to all Platyhelminthes
As bilaterians, platyhelminthes are triploblastic, but have no internal body cavity (are acoelomate), and lack specialized circulatory and respiratory organs, so gas exchange is by simple diffusion. This limits the thickness of the body, so they are either microscopic or are flat and ribbon- or leaf-shaped, and vulnerable to fluid loss. The body is filled with mesenchyme, a connective tissue that can regenerate injured tissues and permits asexual reproduction. The nervous system is concentrated at the head end.
Features specific to Turbellaria
These have about 4,500 species, are mostly free-living, and range from to in length. Most are predators or scavengers, and terrestrial species are mostly nocturnal and live in shaded humid locations such as leaf litter or rotting wood. However some are symbiotes of other animals such as crustaceans, and some are parasites. Free-living turbellarians are mostly black, brown or gray, but some larger ones are brightly colored.
Turbellarians have no cuticle (external layer of organic but non-cellular material). In a few species the skin is a syncitium, a collection of cells with multiple nuclei and a single shared external membrane. However the skins of most species consist of a single layer of cells, each of which generally has multiple cilia (small mobile "hairs"), although in some large species the upper surface has no cilia. These skins are also covered with microvilli between the cilia. They have many glands, usually submerged in the muscle layers below the skin and connect to the surface by pores through which they secrete mucus, adhesives and other substances.
Small aquatic species use the cilia for locomotion, while larger ones use muscular movements of the whole body or of a specialized sole to creep or swim. Some are capable of burrowing, anchoring their rear ends at the bottom of the burrow, then stretching the head up to feed and then pulling it back down for safety. Some terrestrial species throw a thread of mucus which they use as a rope to climb from one leaf to another.
Some Turbelleria have spicular skeletons, giving the appearance of annulations.
Diet and digestion
Most other turbellarians are carnivorous, either preying on small invertebrates or protozoans, or scavenging on dead animals. A few feed on larger animals, including oysters and barnacles, while some, such as Bdelloura, are commensal on the gills of horseshoe crabs. These turbellarians usually have an eversible pharynx, in other words, one that can be extended by being turned inside-out, and the mouths of different species can be anywhere along the underside. The freshwater species Microstomum caudatum can open its mouth almost as wide as its body is long, to swallow prey as large as itself.
The intestine is lined by phagocytic cells which capture food particles that have already been partially digested by enzymes in the gut. Digestion is then completed within the phagocytic cells and the nutrients diffuse through the body.
Nervous system
Concentration of nervous tissue in the head region is least marked in the acoels, which have nerve nets rather like those of cnidarians and ctenophores, but densest around the head. In turbellarians, a distinct brain is present, albeit relatively simple in structure. From the brain one to four pairs of nerve cords run along the length the body, with numerous smaller nerves branching off. The ventral pair of nerve cords are typically the largest, and, in many species, are the only ones present. Unlike more complex animals, such as annelids, there are no ganglia on the nerve cords, other than those forming the brain.
Most turbellarians have pigment-cup ocelli ("little eyes"), one pair in most species, but two or even three pairs in some. A few large species have many eyes in clusters over the brain, mounted on tentacles, or spaced uniformly round the edge of the body. The ocelli can only distinguish the direction from which light is coming and enable the animals to avoid it.
A few groups – mainly catenulids and seriates – have statocysts, fluid-filled chambers containing a small solid particle or, in a few groups, two. These statocysts are thought to be balance and acceleration sensors, as that is the function they perform in cnidarian medusae and in ctenophores. However turbellarian statocysts have no sensory cilia, and it is unknown how they sense the movements and positions of the solid particles.
Most species have ciliated touch-sensor cells scattered over their bodies, especially on tentacles and around the edges. Specialized cells in pits or grooves on the head are probably smell-sensors.
Reproduction
Many turbellarians clone themselves by transverse or longitudinal division, and others, especially acoels, reproduce by budding. The planarian Dugesia is a well-known representative of class Turbellaria.
All turbellarians are simultaneous hermaphrodites, having both female and male reproductive cells, and fertilize eggs internally by copulation. Some of the larger aquatic species mate by penis fencing, a duel in which each tries to impregnate the other, and the loser adopts the female role of developing the eggs.
In turbellarians there are one or more pairs of both testes and ovaries. Sperm ducts run from the testes, through bulb-like seminal vesicles, to the muscular penis. In many species, this basic plan is considerably complicated by the addition of accessory glands or other structures. The penis lies inside a cavity, and can be everted through an opening on the posterior underside of the animal. It often, although not always, possesses a sharp stylet. Unusually among animals, in most species, the sperm cells have two tails, rather than one.
In most species "miniature adults" emerge when the eggs hatch, but a few large species produce plankton-like larvae.
Taxonomy and evolution
Internal relationships
Detailed morphological analyses of anatomical features in the mid-1980s and molecular phylogenetics analyses since 2000 using different sections of DNA agree that Acoelomorpha, consisting of Acoela (traditionally regarded as very simple turbellarians) and Nemertodermatida (another small group previously classified as "turbellarians") are the sister group to all other bilaterians, including the rest of the Platyhelminthes.
The Platyhelminthes is a clade consisting of two monophyletic groups, Catenulida and Rhabditophora.
It has been agreed since 1985 that each of the wholly parasitic platyhelminth groups (Cestoda, Monogenea and Trematoda) is monophyletic, and that together these form a larger monophyletic grouping, the Neodermata, in which the adults of all members have syncitial skins. It is also generally agreed that the Neodermata are a relatively small sub-group a few levels down in the "family tree" of the Rhabditophora. Hence the traditional sub-phylum "Turbellaria" is paraphyletic, since it does not include the Neodermata although these are descendants of a sub-group of "turbellarians".
| Biology and health sciences | Platyzoa | Animals |
662787 | https://en.wikipedia.org/wiki/Elliptic%20partial%20differential%20equation | Elliptic partial differential equation | In mathematics, an elliptic partial differential equation is a type of partial differential equation (PDE). In mathematical modeling, elliptic PDEs are frequently used to model steady states, unlike parabolic PDE and hyperbolic PDE which generally model phenomena that change in time. They are also important in pure mathematics, where they are fundamental to various fields of research such as differential geometry and optimal transport.
Definition
Elliptic differential equations appear in many different contexts and levels of generality.
First consider a second-order linear PDE in two variables, written in the form
where , , , , , , and are functions of and , using subscript notation for the partial derivatives. The PDE is called elliptic if
with this naming convention inspired by the equation for a planar ellipse. Equations with are termed parabolic while those with are hyperbolic.
For a general linear second-order PDE, the "unknown" function can be a function of any number of independent variables; the equation is of the form
where , , and are functions defined on the domain subject to the symmetry . This equation is called elliptic if, when is viewed as a function on the domain valued in the space of symmetric matrices, all of the eigenvalues are greater than some set positive number. Equivalently, this means that there is a positive number such that
for any point in the domain and any real numbers .
The simplest example of a second-order linear elliptic PDE is the Laplace equation, in which is zero if and is one otherwise, and where . The Poisson equation is a slightly more general second-order linear elliptic PDE, in which is not required to vanish. For both of these equations, the ellipticity constant can be taken to be .
The terminology elliptic partial differential equation is not used consistently throughout the literature. What is called "elliptic" by some authors is called strictly elliptic or uniformly elliptic by others.
Nonlinear and higher-order equations
Ellipticity can also be formulated for much more general classes of equations. For the most general second-order PDE, which is of the form
for some given function , ellipticity is defined by linearizing the equation and applying the above linear definition. Since linearization is done at a particular function , this means that ellipticity of a nonlinear second-order PDE depends not only on the equation itself but also on the solutions under consideration. For example, in the simplest kind of Monge–Ampère equation, the determinant of the hessian matrix of a function is prescribed:
As follows from Jacobi's formula for the derivative of a determinant, this equation is elliptic if is a positive function and solutions satisfy the constraint of being uniformly convex.
There are also higher-order elliptic PDE, the simplest example being the fourth-order biharmonic equation. Even more generally, there is an important class of elliptic systems which consist of coupled partial differential equations for multiple 'unknown' functions. For example, the Cauchy–Riemann equations from complex analysis can be viewed as a first-order elliptic system for a pair of two-variable functions.
Moreover, the class of elliptic PDE (of any order, including systems) is subject to various notions of weak solutions, i.e., reformulating the above equations in such a way that allows for solutions to have various irregularities (e.g. non-differentiability, singularities or discontinuities) while still adhering to the laws of physics. Additionally, these type of solutions are also important in variational calculus, where the direct method often produces weak solutions of elliptic systems of Euler equations.
Canonical form
Consider a second-order elliptic partial differential equation
for a two-variable function . This equation is linear in the "leading-order terms" but allows nonlinear expressions involving the function values and their first derivatives; this is sometimes called a quasilinear equation.
A canonical form asks for a transformation and of the domain so that, when is viewed as a function of and , the above equation takes the form
for some new function . The existence of such a transformation can be established locally if , , and are real-analytic functions and, with more elaborate work, even if they are only continuously differentiable. Locality means that the necessary coordinate transformations may fail to be defined on the entire domain of , although they can be established in some small region surrounding any particular point of the domain.
Formally establishing the existence of such transformations uses the existence of solutions to the Beltrami equation. From the perspective of differential geometry, the existence of a canonical form is equivalent to the existence of isothermal coordinates for the associated Riemannian metric
on the domain. (The ellipticity condition for the PDE, namely the positivity of the function , is what ensures that either this tensor or its negation is indeed a Riemannian metric.) Generally, for second-order quasilinear elliptic partial differential equations for functions of more than two variables, a canonical form does not exist. This corresponds to the fact that, although isothermal coordinates generally exist for Riemannian metrics in two dimensions, they only exist for very particular Riemannian metrics in higher dimensions.
Characteristics and regularity
For the general second-order linear PDE, characteristics are defined as the null directions for the associated tensor
called the principal symbol. Using the technology of the wave front set, characteristics are significant in understanding how irregular points of propagate to the solution of the PDE. Informally, the wave front set of a function consists of the points of non-smoothness, in addition to the directions in frequency space causing the lack of smoothness. It is a fundamental fact that the application of a linear differential operator with smooth coefficients can only have the effect of removing points from the wave front set. However, all points of the original wave front set (and possibly more) are recovered by adding back in the (real) characteristic directions of the operator.
In the case of a linear elliptic operator with smooth coefficients, the principal symbol is a Riemannian metric and there are no real characteristic directions. According to the previous paragraph, it follows that the wave front set of a solution coincides exactly with that of . This sets up a basic regularity theorem, which says that if is smooth (so that its wave front set is empty) then the solution is smooth as well. More generally, the points where fails to be smooth coincide with the points where is not smooth. This regularity phenomena is in sharp contrast with, for example, hyperbolic PDE in which discontinuities can form even when all the coefficients of an equation are smooth.
Solutions of elliptic PDEs are naturally associated with time-independent solutions of parabolic PDEs or hyperbolic PDEs. For example, a time-independent solution of the heat equation solves Laplace's equation. That is, if parabolic and hyperbolic PDEs are associated with modeling dynamical systems then the solutions of elliptic PDEs are associated with steady states. Informally, this is reflective of the above regularity theorem, as steady states are generally smoothed out versions of truly dynamical solutions. However, PDE used in modeling are often nonlinear and the above regularity theorem only applies to linear elliptic equations; moreover, the regularity theory for nonlinear elliptic equations is much more subtle, with solutions not always being smooth.
| Mathematics | Differential equations | null |
662807 | https://en.wikipedia.org/wiki/Stun%20grenade | Stun grenade | A stun grenade, also known as a flash grenade, flashbang, thunderflash, or sound bomb, is a non-lethal explosive device used to temporarily disorient an enemy's senses. Upon detonation, a stun grenade produces a blinding flash of light and an extremely loud "bang". They are often used in close-quarters combat, door breaching, and riot control, typically to stun enemies or distract them.
Originally developed to simulate explosions during military training, stun grenades were first used by the British Army Special Air Service's counterterrorist wing in the late 1970s, and have been used by police and military forces worldwide since.
Despite their less-lethal nature, stun grenades are still capable of causing harm, and can injure or kill when detonating in close proximity. They are also capable of sparking fires.
Effects
Stun grenades are designed to produce a blinding flash of light of around 7 megacandela (Mcd) and an intensely loud "bang" of greater than 170 decibels (dB).
Construction
Unlike a fragmentation grenade, stun grenades are constructed with a casing designed to remain intact during detonation and avoid fragmentation injuries, while having large circular cutouts to allow the light and sound of the explosion through. The filler comprises a pyrotechnic metal-oxidant mix of magnesium or aluminium and an oxidizer such as potassium perchlorate or potassium nitrate.
Hazards
While stun grenades are designed to limit injury, permanent hearing loss has been reported. The concussive blast has the ability to cause injuries, and the heat generated may ignite flammable materials. The fires that occurred during the 1980 Iranian Embassy siege in London were caused by stun grenades.
| Technology | Less-lethal weapons | null |
662851 | https://en.wikipedia.org/wiki/Homarus%20gammarus | Homarus gammarus | Homarus gammarus, known as the European lobster or common lobster, is a species of clawed lobster from the eastern Atlantic Ocean, Mediterranean Sea and parts of the Black Sea. It is closely related to the American lobster, H. americanus. It may grow to a length of and a mass of , and bears a conspicuous pair of claws. In life, the lobsters are blue, only becoming "lobster red" on cooking. Mating occurs in the summer, producing eggs which are carried by the females for up to a year before hatching into planktonic larvae. Homarus gammarus is a highly esteemed food, and is widely caught using lobster pots, mostly around the British Isles.
Description
Homarus gammarus is a large crustacean, with a body length up to and weighing up to , although the lobsters caught in lobster pots are usually long and weigh . Like other crustaceans, lobsters have a hard exoskeleton which they must shed in order to grow, in a process called ecdysis (molting). This may occur several times a year for young lobsters, but decreases to once every 1–2 years for larger animals.
The first pair of pereiopods is armed with a large, asymmetrical pair of feet. The larger one is the "crusher", and has rounded nodules used for crushing prey; the other is the "cutter", which has sharp inner edges, and is used for holding or tearing the prey. Usually, the left claw is the crusher, and the right is the cutter.
The exoskeleton is generally blue above, with spots that coalesce, and yellow below. The red colour associated with lobsters only appears after cooking. This occurs because, in life, the red pigment astaxanthin is bound to a protein complex, but the complex is broken up by the heat of cooking, releasing the red pigment.
The closest relative of H. gammarus is the American lobster, Homarus americanus. The two species are very similar, and can be crossed artificially, although hybrids are unlikely to occur in the wild since their ranges do not overlap. The two species can be distinguished by a number of characteristics:
The rostrum of H. americanus bears one or more spines on the underside, which are lacking in H. gammarus.
The spines on the claws of H. americanus are red or red-tipped, while those of H. gammarus are white or white-tipped.
The underside of the claw of H. americanus is orange or red, while that of H. gammarus is creamy white or very pale red.
Life cycle
Female H. gammarus reach sexual maturity when they have grown to a carapace length of , whereas males mature at a slightly smaller size. Mating typically occurs in summer between a recently moulted female, whose shell is therefore soft, and a hard-shelled male. The female carries the eggs attached to her pleopods for up to 12 months, depending on the temperature. Females carrying eggs are said to be "berried" and can be found throughout the year.
The eggs hatch at night, and the larvae swim to the water surface where they drift with the ocean currents, preying on zooplankton. This stage involves three moults and lasts for 15–35 days. After the third moult, the juvenile takes on a form closer to the adult, and adopts a benthic lifestyle. The juveniles are rarely seen in the wild, and are poorly known, although they are known to be capable of digging extensive burrows. It is estimated that only 1 larva in every 20,000 survives to the benthic phase. When they reach a carapace length of , the juveniles leave their burrows and start their adult lives.
Distribution
Homarus gammarus is found across the north-eastern Atlantic Ocean from northern Norway to the Azores and Morocco, not including the Baltic Sea. It is also present in most of the Mediterranean Sea, only missing from the section east of Crete, and along only the south-west coast of the Black Sea. The northernmost populations are found in the Norwegian fjords Tysfjorden and Nordfolda, inside the Arctic Circle.
The species can be divided into four genetically distinct populations, one widespread population, and three which have diverged due to small effective population sizes, possibly due to adaptation to the local environment. The first of these is the population of lobsters from northern Norway, which have been referred to as the "midnight-sun lobster". The populations in the Mediterranean Sea are distinct from those in the Atlantic Ocean. The last distinct population is found in part of the Netherlands: samples from the Oosterschelde were distinct from those collected in the North Sea or English Channel.
Attempts have been made to introduce H. gammarus to New Zealand, alongside other European species such as the edible crab, Cancer pagurus. Between 1904 and 1914, one million lobster larvae were released from hatcheries in Dunedin, but the species did not become established there.
Ecology
Adult H. gammarus live on the continental shelf at depths of , although not normally deeper than . They prefer hard substrates, such as rocks or hard mud, and live in holes or crevices, emerging at night to feed.
The diet of H. gammarus mostly consists of other benthic invertebrates. These include crabs, molluscs, sea urchins, starfish and polychaete worms.
The three clawed lobster species Homarus gammarus, H. americanus and Nephrops norvegicus are hosts to the three known species of the animal phylum Cycliophora; the species on H. gammarus has not been described.
Homarus gammarus is susceptible to the disease gaffkaemia, caused by the bacterium Aerococcus viridans. Although it is frequently found in American lobsters, the disease has only been seen in captive H. gammarus, where prior occupation of the tanks by H. americanus could not be ruled out.
Human consumption
Homarus gammarus is traditionally "highly esteemed" as a foodstuff and was mentioned in "The Crabfish" a seventeenth century English folk song. It may fetch very high prices and may be sold fresh, frozen, canned or powdered. Both the claws and the abdomen of H. gammarus contain "excellent" white meat, and most of the contents of the cephalothorax are edible. The exceptions are the gastric mill and the "sand vein" (gut). The price of H. gammarus is up to three times higher than that of H. americanus, and the European species is considered to be more flavorful.
Lobsters are mostly fished using lobster pots, although lines baited with octopus or cuttlefish sometimes succeed in tempting them out, to allow them to be caught in a net or by hand. In 2008, 4,386 t of H. gammarus were caught across Europe and North Africa, of which 3,462 t (79%) was caught in the British Isles (including the Channel Islands). The minimum landing size for H. gammarus is a carapace length of . To protect known breeding females, lobsters caught carrying eggs are to be notched on a uropod, the inner tail flap of female lobsters of reproductive size (usually above the minimum landing size 87mm carapace length). Following this, it is illegal for the female to be kept or sold, and is commonly referred to as a "v-notch". This notch remains for three molts of the lobster exoskeleton, providing harvest protection and continued breeding availability for 3–5 years.
Aquaculture systems for H. gammarus are under development, and production rates are still very low.
Taxonomic history
Homarus gammarus was first given a binomial name by Carl Linnaeus in the tenth edition of his Systema Naturae, published in 1758. That name was Cancer gammarus, since Linnaeus' concept of the genus Cancer at that time included all large crustaceans.
H. gammarus is the type species of the genus Homarus Weber, 1795, as determined by Direction 51 of the International Commission on Zoological Nomenclature. Prior to that direction, confusion arose because the species had been referred to by several different names, including Astacus marinus Fabricius, 1775 and Homarus vulgaris H. Milne-Edwards, 1837, and also because Friedrich Weber's description of the genus had been overlooked until rediscovered by Mary J. Rathbun, rendering any prior assignments of type species (for Homarus H. Milne-Edwards, 1837) invalid for Homarus Weber, 1795.
The type specimen of Homarus gammarus was a lectotype selected by Lipke Holthuis in 1974. It came from , near Marstrand, Sweden ( northwest of Gothenburg), but both it and the paralectotypes have since been lost.
The common name for H. gammarus preferred by the Food and Agriculture Organization is "European lobster", but the species is also widely known as the "common lobster".
| Biology and health sciences | Crayfishes and lobsters | Animals |
662879 | https://en.wikipedia.org/wiki/Mobile%20device | Mobile device | A mobile device or handheld device is a computer small enough to hold and operate in hand. Mobile devices are typically battery-powered and possess a flat-panel display and one or more built-in input devices, such as a touchscreen or keypad. Modern mobile devices often emphasize wireless networking, to both the Internet and to other devices in their vicinity, such as headsets or in-car entertainment systems, via Wi-Fi, Bluetooth, cellular networks, or near-field communication.
Characteristics
Device mobility can be viewed in the context of several qualities:
Physical dimensions and weight
Whether the device is mobile or some kind of host to which it is attached is mobile
What kind of host devices it can be bound with
How devices communicate with a host
When mobility occurs
Strictly speaking, many so-called mobile devices are not mobile. It is the host that is mobile, i.e., a mobile human host carries a non-mobile smartphone device. An example of a true mobile computing device, where the device itself is mobile, is a robot. Another example is an autonomous vehicle.
There are three basic ways mobile devices can be physically bound to mobile hosts:
Accompanied,
Surface-mounted, or
Embedded into the fabric of a host, e.g., an embedded controller in a host device.
Accompanied refers to an object being loosely bound and accompanying a mobile host, e.g., a smartphone can be carried in a bag or pocket but can easily be misplaced. Hence, mobile hosts with embedded devices such as an autonomous vehicle can appear larger than pocket-sized.
The most common size of a mobile computing device is pocket-sized, but other sizes for mobile devices exist. Mark Weiser, known as the father of ubiquitous computing, referred to device sizes that are tab-sized, pad, and board sized, where tabs are defined as accompanied or wearable centimeter-sized devices, e.g. smartphones, phablets and tablets are defined as hand-held decimeter-sized devices. If one changes the form of the mobile devices in terms of being non-planar, one can also have skin devices and tiny dust-sized devices.
Dust refers to miniaturized devices without direct HCI interfaces, e.g., micro-electromechanical systems (MEMS), ranging from nanometers through micrometers to millimeters. | Technology | Computer hardware | null |
662948 | https://en.wikipedia.org/wiki/Auckland%20Harbour%20Bridge | Auckland Harbour Bridge | The Auckland Harbour Bridge is an eight-lane motorway bridge over the Waitematā Harbour in Auckland, New Zealand. It joins St Marys Bay on the Auckland city side with Northcote on the North Shore side. It is part of State Highway 1 and the Auckland Northern Motorway. The bridge is operated by the NZ Transport Agency (NZTA). It is the second-longest road bridge in New Zealand, and the longest in the North Island.
The original inner four lanes, opened in 1959, are of box truss construction. Two lanes were added to each side in 1968–1969 and are of orthotropic box structure construction extend as cantilevers from the original piers. The bridge is 1,020 m (3,348 ft) long, with a main span of 243.8 metres (800 feet) rising 43.27 metres (142 feet) above high water, allowing ships access to the deepwater wharf at the Chelsea Sugar Refinery, one of the few such wharves west of the bridge.
While often considered an Auckland icon, many see the construction of the bridge without walking, cycling, and rail facilities as a big oversight. In 2016, an add-on structure providing a walk-and-cycleway called SkyPath received Council funding approval and planning consent, but was not built. In 2021, a stand-alone walking and cycling bridge called the Northern Pathway was announced by the New Zealand Government, but also was not built.
About 170,000 vehicles cross the bridge each day (as of 2019), including over 1,000 buses, which carry 38% of all people crossing during the morning peak.
Background
Prior to the opening of the bridge in 1959, the quickest way from Auckland to the North Shore was by passenger or vehicular ferry. By road, the shortest route was via the Northwestern Motorway (then complete only between Great North Road and Lincoln Road), Massey, Riverhead, and Albany, a distance of approximately .
As early as 1860, engineer Fred Bell, commissioned by North Shore farmers who wanted to herd animals to market in Auckland, had proposed a harbour crossing in the general vicinity of the bridge. It would have used floating pontoons, but the plan failed due to the £16,000 cost estimate ($1.9 million, adjusted for inflation as of March 2017). Additional structures for a bridge crossing the harbour were proposed in 1927 and 1929.
In the 1950s, when the bridge was being built, North Shore was a mostly rural area of barely 50,000 people, with few jobs and a growth rate half that of Auckland south of the Waitematā Harbour. Opening up the area via a new route unlocked the potential for further expansion of Auckland.
Construction
Initial structure
The recommendations of the design team and the report of the 1946 Royal Commission were for five or six traffic lanes, with one or two of them to be reversed in direction depending on the flow of traffic, and with a footpath for pedestrians on each side. The latter features were dropped for cost reasons before construction started, the First National Government of New Zealand opting for an 'austerity' design of four lanes without footpaths, and including an approach road network only after local outcry over traffic effects. The decision to reduce the bridge in this way has been called "a ringing testament to [...] the peril of short-term thinking and penny-pinching". On 1 December 1950, an act of parliament formed the Auckland Harbour Bridge Authority, chaired by Sir John Allum, then Mayor of Auckland City, who appointed British firm Freeman Fox & Partners to design the bridge.
The bridge took four years to build, with Dorman Long (who had constructed the Sydney Harbour Bridge) and the Cleveland Bridge & Engineering Company contracted to construct the bridge in October 1954. The first stage of construction involved land reclamation at the Westhaven Marina, which was completed by September 1955. The steel girder structure pieces were fabricated in England and shipped to New Zealand. The steel bridge structure began erection in December 1956. Hundreds of labourers were employed on the construction including 180 men sent out from the UK. Progress was slowed with the workers going on strike in 1956 and 1957. The large steel girder sections were partially pre-assembled, then floated into place on construction barges. One of the main spans was almost lost during stormy weather when the barge began to drift, but the tugboat William C Daldy won a 36-hour tug-of-war against the high winds.
The bridge was constructed from opposing sides of the harbour. The southern section was cantilevered, until both sides were joined in March 1959. Completed in April 1959, three weeks ahead of schedule, the bridge was officially opened on 30 May 1959 by the Governor-General Lord Cobham. An open day had been held, when 106,000 people had walked across. The opening period was extremely busy, despite the poor weather in Auckland experienced in June 1959. Either three or four men had been killed by accidents during construction, and the names of three of them are recorded on a memorial plaque underneath the bridge at the Northcote end.
The hollow girder design by Freeman, Fox and Partners design was unprecedented in New Zealand, and fell outside the 1950s building codes in New Zealand. Initial plans for the bridge were for an extremely slender structure, only 2.9 metres thick, due to the competing specifications from two stakeholders: the National Roads Board specified the gradient and locations where the bridge could launch from the shore on either side of the harbour, while the Auckland Harbour Board required an opening of 43.5 metres above the high tide point. Public Works commissioner Bob Norman, concerned about the narrow bridge design, attempted to negotiate with both the Roads Board and Harbour Board for additional width allowance for the bridge. The Harbour Board required the 43.5 metre clearance so that the entire fleet of ships operating within New Zealand could navigate the harbour, the largest of which was the P&O cruise liner SS Canberra. Norman argued that the Canberra was extremely unlikely to use the only major dock west of the bridge at the Chelsea Sugar Refinery, so the Harbour Board agreed to a smaller opening. This allowed Freeman Fox and Partners to redesign the bridge, increasing the width of the deep centre span from 2.9 metres to 4.12 metres. By the 1970s, many box girder bridges began to develop structural problems, such as the Freeman Fox and Partners-designed West Gate Bridge in Melbourne which collapsed during construction in 1970. The Auckland Harbour Bridge was inspected by the design firm, which found that the stiffening member had buckled by 61mm, so it was decided to strengthen the bridge's girder system.
Paid for by government-backed loans, the bridge started out as a toll bridge, the first one in New Zealand, with toll booths at the northern end for north-bound and south-bound traffic. Tolls were originally 2/6 (2 shillings and six pence: approximately $5.50 in 2018) per car but were reduced to 2/- (2 shillings: approximately $4.47 in 2018) after 15 months of operation. The toll remained at 2 shillings until New Zealand changed to decimal currency in July 1967, when that amount became 20 cents in the conversion. It was increased in 1980 from 20 to 25 cents (approximately $1.21 in 2018). Tolling was later made north-bound only before being discontinued on 31 March 1984, and the booths were removed. The toll system was removed as the cost of collection began to outweigh the profits. When this happened, the Auckland Harbour Bridge Authority enquired if the National Roads Board would take over operations if the toll booths were removed, which they agreed to. When the bridge became toll free, most of the Auckland Harbour Bridge Authority staff were absorbed into the roads board.
Some critics have alleged that the routing of State Highway 1 over the bridge was motivated by the need to create toll revenue, and led to a decades-long delay on finishing the Western Ring Route around Auckland, significantly contributing to the need for a massive motorway through the city centre of Auckland and severely damaging inner-city suburbs such as Freemans Bay and Grafton.
Clip-ons (Nippon clip-ons)
The bridge was originally built with four lanes for traffic. Owing to the rapid expansion of suburbs on the North Shore and increasing traffic levels, it was soon necessary to increase capacity; by 1965, the annual use was about 10 million vehicles, three times the original forecast. In 1967, a contract was given to Japanese firm Ishikawajima-Harima Heavy Industries Co., Ltd. (now IHI Corporation) to construct two steel box girder bridges affixed to the Harbour Bridge, to greatly increase the number of lanes on the bridge. The girder sections were prefabricated in Japan and transported to New Zealand on a converted oil tanker. The eastern section was completed in January 1969, while the western side was completed shortly before the additional lanes were formally opened on 23 September 1969. Each side added two additional lanes to the bridge, doubling the number of lanes to eight.
As the sections were manufactured by a Japanese company, this led to the nickname 'Nippon clip-ons'. The selection of the company was considered a bold move at the time, barely 20 years after WWII and with some considerable anti-Japanese sentiment still existing. The costs of the additions were much higher than if the extra lanes had been provided initially.
The clip-ons have been plagued by significant issues. In 1987, cracks required major repair works, and in 2006, further cracks and signs of material fatigue were found. The clip-ons were originally to have a life expectancy of 50 years. Auckland City Council's Transport Committee requested Transit New Zealand to investigate the future of the clip-ons as part of its ten-year plan. Transit noted that the plan already includes some funding for bridge maintenance.
In May 2007, Transit proposed a by-law change banning vehicles over 4.5 tonnes from the outside lane on each clip-on to reduce stress on the structure. This was changed in July 2007 to a bylaw banning vehicles of 13 tonnes or more, based on the high level of voluntary compliance during the previous months.
In 2007, it was announced that NZ$45 million in maintenance work on the clip-ons was brought forward as part of good practice. In October 2007, a 2006 report from Beca Group surfaced in the press, noting that the clip-ons were at risk of catastrophic, immediate failure in circumstances such as a traffic jam trapping a large number of trucks. Transit noted that this situation was extremely unlikely, and measures already implemented would prevent it from occurring. In January 2008, it became known that even after the multimillion-dollar maintenance works, a full ban for trucks on all clip-on lanes might be required, or the working life could be reduced to only ten more years.
In late 2009, it was announced that due to greater than expected complexity of the task and increasing material costs for the 920 tons of reinforcing material instead of the approximately half amount of that originally envisaged, clip-on maintenance costs had increased by a further NZ$41 million. NZTA noted that the clip-ons would not be able to be strengthened again after the current works were finished. However, after completion of the upgrade, the bridge would have a further life of between 20 and 40 years if truck restrictions were reintroduced in 10–20 years on the northbound clip-on.
Traffic
Tidal flow
A "tidal flow" (dynamic lanes) system is in place, with the direction of the two centre lanes changed to provide an additional lane for peak-period traffic. During the morning peak, five of the eight lanes are for southbound traffic; in the afternoon, five lanes are northbound. At other times, the lanes are split evenly, but peak traffic has become proportionately less – in 1991 there was often a higher than 3:1 difference in directional traffic; in 2006, this had dropped to around 1.6:1. The bridge has an estimated capacity of 180,000 vehicles per day, and in 2006 had an average volume of 168,754 vehicles per day (up from 122,000 in 1991).
In March 1982, the Ministry of Transport and Auckland Harbour Bridge Authority conducted a week-long traffic blitz in an attempt to improve the standard of driving. Of the 600,000 vehicles which used the bridge over this period, 6,000 were stopped, with half of those receiving a ticket and the rest cautioned. A second blitz was held for 36 hours a few weeks later.
For many years, lane directions were indicated by overhead signals. In the late 1980s, a number of fatal head-on accidents occurred when vehicles crossed lane markings into the path of oncoming traffic.
Moveable barriers
In 1990, a movable concrete safety barrier was put in place to separate traffic heading in opposite directions and eliminate head-on accident. Two specially designed barrier transfer machines moved the barrier by one lane four times a day, at a speed of 6 km/h, the first concrete safety barrier of its kind installed on a box girder bridge in the world.
In March 2009, the barrier transfer machines, which had lasted four times their original design life of five years, and the barrier were replaced. The new machines are capable of moving the barrier in half the time the old machines did. The concrete barrier blocks and the metal expansion blocks have been reduced in width by 200 mm, giving more width in the lanes either side of the barrier.
As part of the Victoria Park Tunnel project, the moveable barrier has been extended southwards to the Fanshawe Street onramp.
Event management
As part of large events such as the Auckland Marathon, normal motorway restrictions on access are sometimes relaxed. December 2011 was the first time that cyclists were officially allowed on the bridge, for a race / community cycling event organised by Telstra Clear, Auckland Transport, NZTA and Cycle Action Auckland, also allowing cyclists on the Northern Busway. Up to 9,000 riders were protected by 160 stationary buses used as a 'guard of honour' between the bridge end and the Northern Busway from traffic on the rest of the motorway.
Proposed walk- and cycleway
Original proposal
When the bridge was built, rail lines and walking paths were dropped for cost reasons, and neither were they included during the clip-on construction (people can walk on the span only via guided tours). After the early 1990s increase in public transportation patronage in Auckland, the Ministry of Works and Development investigated if the 'clip-ons' could be used for a light rail system, which they found was feasible if the lanes were used exclusively for this purpose. In 2007 discussions about the addition of a cycle and footpath link were mooted. Transit noted that this would cost between NZ$20 million and $40 million, but public support was polled as very high. The GetAcross group and Cycle Action Auckland (later rebranded to Bike Auckland) argued that lower-cost options were available, and that provision for a walk- and cycleway could relatively easily be included in the bridge strengthening works that were being planned for the clip-ons.
A 2008 proposal to modify the clip-ons and potentially widen them to add walking and cycling paths met with different reactions. While Auckland Regional Council and North Shore City Council voted to support it (under certain conditions), Auckland City Council considered the costs to be too high. Other stakeholders such as the NZ Transport Agency (NZTA) considered the proposal as not having enough merit for the $22–53 million cost, though campaigners noted that the costs cited for the project included 45% contingencies. A proposal from the Auckland Regional Council (one of the proponents) to open up part of the clip-on structure for a walking / cycling trial use over several summer weekends, to show whether it would attract enough users, did not go forward.
The GetAcross group was showcasing its proposed walking/cycling solution, called SkyPath, on its website. Following years of campaigning a Harbour Bridge crossing, known as Skypath, was promised funding by the Labour Party in the lead-up to the 2017 general election. Once Labour was in government, the project was passed to the Waka Kotahi / NZ Transport Agency which released a revised design in 2019.
In July 2023 Bike Auckland released a report by SmartSense Limited, addressing key concerns about reallocating a lane on the motor bridge to walking and cycling, and proposing a design solution to mitigate safety concerns.
On 6 August 2023, Waka Kotahi announced their Waitematā Harbour Crossings plan which includes a tunnel for light rail and a tunnel for motor traffic under the Harbour, and walking and cycling on two lanes of the existing Harbour Bridge. Construction is expected to start by 2029. Waka Kotahi's forecast is that 6400 people would walk and cycle across the Auckland Harbour Bridge every day.
Bike Auckland continues to advocate for Waka Kotahi to Liberate the Lane, stating that Waka Kotahi's Waitematā connections project will take too long to deliver a walking and cycling connection across the Harbour. Their campaign has attracted the support of a diverse array of organisations, calling for Waka Kotahi to liberate the lane now to give Aucklanders more affordable and sustainable transport options, and that it would be a key symbol of climate action.
Protests
On Sunday, 24 May 2009, thousands of people crossed the bridge as a part of a protest by GetAcross against the bridge not providing walking and cycling access, and against what the group perceives to be the authorities' negative and obstructionist attitude towards such access. A crossing either as part of the protest or as part of the official 50-year anniversary celebrations had been forbidden by NZTA because of the costs and traffic difficulties claimed for a managed crossing. However, after several speeches, including by Auckland Regional Council Chairman Mike Lee, several people made their way around the police cordon onto the bridge. At that stage police closed the northbound lanes to traffic, bringing State Highway 1 to a stop. The remainder of the protesters moved onto the bridge, which was not resisted any more by the police. No accidents, violence or arrests were reported, and protesters left the bridge approximately an hour later, many having crossed to the North Shore and back.
The protest created a wide spectrum of responses in the media and in public perception, from being labelled a dangerous stunt representative of an increasingly lawless, anarchic society to being considered a successful signal to authorities to give more weight to the demands and the public backing of the walk and cycleway proponents. Authorities noted that they were investigating whether any of the protesters would face fines or charges. NZTA representatives noted that they were disappointed at what they considered the broken word of the organisers of the protest, and remarked that it would take 30 more years before walking and cycling could likely be provided (see also "Second Harbour Crossing" below). NZTA were criticised as having brought the situation at least partly onto themselves by choosing the easy route of forbidding the protest crossing. Several political protest marches (especially hīkoi) had been allowed to cross the bridge.
Updated proposal
Because of the costs of the proposal and increasing information about the problematic state of the clip-ons, the GetAcross campaign in late 2009 proposed an alternative solution, with a single shared walking and cycling path slung under the eastern clip-on. As confirmed by NZTA, this clip-on has significantly more remaining load capacity (it is used by fewer heavy trucks, being the route of (often empty) trucks returning to Ports of Auckland) and as the proposal would not require widening, the costs have been preliminarily assessed as of the order of NZ$12 million.
The group proposes to raise the majority of the funding via a loan backed by small tolls, of the order of NZ$1 for regular users. NZTA noted that it would be considering the proposal, should funding be able to be secured by the campaigners.
In 2011, the proposal got new public support when Auckland Mayor Len Brown agreed that a walk- and cycleway was a desirable goal, and instructed Auckland Transport to add it to its strategic priorities. The walk- and cycleway is also to be included in the city centre masterplan. Three council-controlled organisations (CCOs) – Auckland Transport, the Waterfront Development Agency and the Tourism, Events and Economic Development Agency – indicated support for the proposal, as has the Heart of the City (Auckland CBD) business association.
In August 2011, an editorial in The New Zealand Herald gave conditional support to the newest proposal, noting that a toll-based funding model and the partially enclosed weather-protected design of the $23 million proposal by Hopper Developments would appear to cover most concerns.
Approval
In 2014, the proposed walk and cycleway was publicly notified, and consent was given in 2015. However, this was appealed by three local groups (two which later dropped out of the appeal). The decision of the original hearing was upheld in December 2016, and the last appeal rejected by the Environment Court. In the meantime, Council had already provided in principle approval for a public-private partnership funding model, in a unanimous support vote earlier in 2016.
A 2019 announcement said that work on the walking and cycling "clipon" could start in 2020. Mayoral candidate John Tamihere proposed replacement with a 10-lane lower level plus rail and cycling/pedestrian facilities on an upper level.
Protests over lack of progress
On 30 May 2021, more than 1,500 cyclists crossed the bridge following the Liberate the Lane rally at Point Erin Park organised by Bike Auckland. The rally called for a trial of reallocating a traffic lane for walking and cycling on the bridge over the summer and included speeches by Auckland Central MP Chlöe Swarbrick and former associate minister of transport Julie Anne Genter. The rally was motivated by uncertainty around the future of the SkyPath project. Waka Kotahi had quietly sidelined the project due to technical issues. The Western clip on of the bridge (two motor traffic lanes) had been closed in advance of the rally, with a police cordon blocking access. After Bike Auckland's rally concluded, much of the crowd made their way over to the police cordon and pushed past onto the bridge, to show their determination for access for walking and cycling to be provided on the Auckland Harbour Bridge. No injuries were reported however one person was arrested for breaching the cordon, before being released without charge.
In response to the proposal for a trial cycle lane, NZTA stated that a cycle lane would likely require two lanes in order to provide sufficient protection for cyclists and pedestrians.
Stand-alone bridge
A couple of days after Bike Auckland's rally, in June 2021, Transport Minister Michael Wood announced a new stand-alone walking and cycling bridge would be built on the eastern side of the Auckland Harbour Bridge. The bridge was estimated to cost a total of $785 million and had the support of Auckland mayor Phil Goff who said it would benefit both Aucklanders and tourists.
The plan received criticism from cycling, trucking and other transport advocates, as well as from the government opposition parties. In October 2021, Wood announced the project had been scrapped due to lack of public support. He said Waka Kotahi had spent $51 million on designs, consultants and engineering plans for the project up until the end of September, and the final amount spent was not known.
Liberate the lane
In 2022 Waka Kotahi confirmed it would not provide a trial of walking and cycling on the Auckland Harbour Bridge due to concerns around safety of people using the lane and motor congestion on the bridge.
In July 2023 Bike Auckland released a report by consultant SmartSense Limited, addressing Waka Kotahi's key concerns about reallocating a lane on the motor bridge to walking and cycling, including a design solution to mitigate safety concerns. The report revealed that motor traffic volumes have declined, leaving space on the bridge to reallocate one lane for walking, cycling, and wheeling "without significantly affecting motor traffic" Bike Auckland continues to advocate to Waka Kotahi to Liberate the Lane, stating that Waka Kotahi's Waitematā connections project will take too long to deliver a walking and cycling connection across the Harbour. Their campaign has attracted the support of a diverse array of organisations, all calling for Waka Kotahi to liberate the lane now. Reasons for their support range from giving Aucklanders more affordable and sustainable transport options, to it being a key action for climate action mitigation and emissions reduction.
Utilities
The bridge supports several utility services, including water and gas pipelines and fibre-optic telecommunications cables.
Transpower reached agreement with Transit in 2005 for the installation of cable supports beneath the bridge for a future cross-harbour power cable. In 2012, Transpower installed three 220,000-volt cables on the bridge, linking Hobson Street substation in the Auckland CBD to the Wairau Road substation on the North Shore.
Tourism
Bungy jumping and Bridge Climb
AJ Hackett operates a bungy jump experience and a guided bridge climb over the arch truss. In popular culture, Bryan Bruce's television documentary The Bridge (2002) featured footage of the first bungy jump from the Auckland Harbour Bridge.
Vector lights
Vector Limited, a utility company in New Zealand, uses LED lights with various colours to illuminate the bridge for ornamental reasons. The lights are powered by renewable energy and the installation was completed in November 2017.
Second Harbour Crossing
Almost as soon as the bridge was built it reached capacity, before extension via the clip-ons, and a second crossing of the harbour was mooted. The high costs and the difficulties of connecting it to the motorway network have so far caused plans to remain at concept stage. In 2008, a study group narrowed down around 160 options to a multi-tunnel link approximately one km east of the bridge, with up to four individual tunnels for motorway, public transport and rail.
On 6 August 2023, Waka Kotahi announced their Waitematā Harbour Crossings plan which includes a tunnel for light rail and a tunnel for motor traffic, with walking and cycling on two lanes of the existing Harbour Bridge. Construction is expected to start by 2029. Waka Kotahi's forecast is that 6,400 people would walk and cycle across the bridge every day.
Issues and notable incidents
Resonance
The natural sway motion of people walking on the bridge's clip-on segment during special events such as the Auckland Marathon can lead the bridge to oscillate sideways. It has been reported that the oscillations can inflict "serious crushing injuries". The bridge's movement is caused by synchronous lateral excitation, a positive feedback phenomenon. It has been a known issue since at least 1975, and the lateral frequency is reported to be at 0.67 Hz.
Seismic vulnerability and improvements
In 1996, NZTA began a seismic screening programme to identify existing bridges that may sustain damage in an earthquake. Following the assessment, the bridge has completed seismic retrofit.
Suicides
The bridge is associated with suicide attempts. In 2010, a news article reported that one to two individuals die by suicide at the location per year. In 2019, a feasibility study of retrofitting suicide prevention barrier was published. It examined two types of barriers, a vertical anti-climb barrier and horizontal fall prevention barrier. As of 2024, no barriers have been implemented.
2020 structural damage
On Friday, 18 September 2020, at approximately 11:00 AM, high winds, with gusts up to , caused a heavy goods vehicle travelling in the central span of the bridge to strike a main diagonal member of the box truss member. The incident reportedly caused severe damage and a significant reduction in daily traffic capacity. Temporary repairs were effected using a locally fabricated replacement member, pending a full engineering analysis and design of the long-term solution. On 4 October, a permanent replacement strut was installed, with all lanes opening again on 7 October 2020. Since the September 2020 incident up until 2024, the bridge has been closed 20 times totalling 435 minutes and costing over $35.6 million in economic damage.
Ship strike risks
The Baltimore Key Bridge collapse in March 2024 prompted NZ Transport Agency to review Auckland Harbour Bridge safety measures. The review classified the risk of vessel collision with the bridge as "rare," citing multiple existing safety measures. The review also stated that Chelsea sugar ship is the only cargo ship that regularly passes under the bridge.
| Technology | Bridges | null |
663041 | https://en.wikipedia.org/wiki/Greatest%20element%20and%20least%20element | Greatest element and least element | In mathematics, especially in order theory, the greatest element of a subset of a partially ordered set (poset) is an element of that is greater than every other element of . The term least element is defined dually, that is, it is an element of that is smaller than every other element of
Definitions
Let be a preordered set and let
An element is said to be if and if it also satisfies:
for all
By switching the side of the relation that is on in the above definition, the definition of a least element of is obtained. Explicitly, an element is said to be if and if it also satisfies:
for all
If is also a partially ordered set then can have at most one greatest element and it can have at most one least element. Whenever a greatest element of exists and is unique then this element is called greatest element of . The terminology least element of is defined similarly.
If has a greatest element (resp. a least element) then this element is also called (resp. ) of
Relationship to upper/lower bounds
Greatest elements are closely related to upper bounds.
Let be a preordered set and let
An is an element such that and for all Importantly, an upper bound of in is required to be an element of
If then is a greatest element of if and only if is an upper bound of in In particular, any greatest element of is also an upper bound of (in ) but an upper bound of in is a greatest element of if and only if it to
In the particular case where the definition of " is an upper bound of " becomes: is an element such that and for all which is to the definition of a greatest element given before.
Thus is a greatest element of if and only if is an upper bound of .
If is an upper bound of that is not an upper bound of (which can happen if and only if ) then can be a greatest element of (however, it may be possible that some other element a greatest element of ).
In particular, it is possible for to simultaneously have a greatest element for there to exist some upper bound of .
Even if a set has some upper bounds, it need not have a greatest element, as shown by the example of the negative real numbers.
This example also demonstrates that the existence of a least upper bound (the number 0 in this case) does not imply the existence of a greatest element either.
Contrast to maximal elements and local/absolute maximums
A greatest element of a subset of a preordered set should not be confused with a maximal element of the set, which are elements that are not strictly smaller than any other element in the set.
Let be a preordered set and let
An element is said to be a if the following condition is satisfied:
whenever satisfies then necessarily
If is a partially ordered set then is a maximal element of if and only if there does exist any such that and
A is defined to mean a maximal element of the subset
A set can have several maximal elements without having a greatest element.
Like upper bounds and maximal elements, greatest elements may fail to exist.
In a totally ordered set the maximal element and the greatest element coincide; and it is also called maximum; in the case of function values it is also called the absolute maximum, to avoid confusion with a local maximum.
The dual terms are minimum and absolute minimum.
Together they are called the absolute extrema.
Similar conclusions hold for least elements.
Role of (in)comparability in distinguishing greatest vs. maximal elements
One of the most important differences between a greatest element and a maximal element of a preordered set has to do with what elements they are comparable to.
Two elements are said to be if or ; they are called if they are not comparable.
Because preorders are reflexive (which means that is true for all elements ), every element is always comparable to itself.
Consequently, the only pairs of elements that could possibly be incomparable are pairs.
In general, however, preordered sets (and even directed partially ordered sets) may have elements that are incomparable.
By definition, an element is a greatest element of if for every ; so by its very definition, a greatest element of must, in particular, be comparable to element in
This is not required of maximal elements.
Maximal elements of are required to be comparable to every element in
This is because unlike the definition of "greatest element", the definition of "maximal element" includes an important statement.
The defining condition for to be a maximal element of can be reworded as:
For all (so elements that are incomparable to are ignored) then
Example where all elements are maximal but none are greatest
Suppose that is a set containing (distinct) elements and define a partial order on by declaring that if and only if
If belong to then neither nor holds, which shows that all pairs of distinct (i.e. non-equal) elements in are comparable.
Consequently, can not possibly have a greatest element (because a greatest element of would, in particular, have to be comparable to element of but has no such element).
However, element is a maximal element of because there is exactly one element in that is both comparable to and that element being itself (which of course, is ).
In contrast, if a preordered set does happen to have a greatest element then will necessarily be a maximal element of and moreover, as a consequence of the greatest element being comparable to element of if is also partially ordered then it is possible to conclude that is the maximal element of
However, the uniqueness conclusion is no longer guaranteed if the preordered set is also partially ordered.
For example, suppose that is a non-empty set and define a preorder on by declaring that holds for all The directed preordered set is partially ordered if and only if has exactly one element. All pairs of elements from are comparable and element of is a greatest element (and thus also a maximal element) of So in particular, if has at least two elements then has multiple greatest elements.
Properties
Throughout, let be a partially ordered set and let
A set can have at most greatest element. Thus if a set has a greatest element then it is necessarily unique.
If it exists, then the greatest element of is an upper bound of that is also contained in
If is the greatest element of then is also a maximal element of and moreover, any other maximal element of will necessarily be equal to
Thus if a set has several maximal elements then it cannot have a greatest element.
If satisfies the ascending chain condition, a subset of has a greatest element if, and only if, it has one maximal element.
When the restriction of to is a total order ( in the topmost picture is an example), then the notions of maximal element and greatest element coincide.
However, this is not a necessary condition for whenever has a greatest element, the notions coincide, too, as stated above.
If the notions of maximal element and greatest element coincide on every two-element subset of then is a total order on
Sufficient conditions
A finite chain always has a greatest and a least element.
Top and bottom
The least and greatest element of the whole partially ordered set play a special role and are also called bottom (⊥) and top (⊤), or zero (0) and unit (1), respectively.
If both exist, the poset is called a bounded poset.
The notation of 0 and 1 is used preferably when the poset is a complemented lattice, and when no confusion is likely, i.e. when one is not talking about partial orders of numbers that already contain elements 0 and 1 different from bottom and top.
The existence of least and greatest elements is a special completeness property of a partial order.
Further introductory information is found in the article on order theory.
Examples
The subset of integers has no upper bound in the set of real numbers.
Let the relation on be given by The set has upper bounds and but no least upper bound, and no greatest element (cf. picture).
In the rational numbers, the set of numbers with their square less than 2 has upper bounds but no greatest element and no least upper bound.
In the set of numbers less than 1 has a least upper bound, viz. 1, but no greatest element.
In the set of numbers less than or equal to 1 has a greatest element, viz. 1, which is also its least upper bound.
In with the product order, the set of pairs with has no upper bound.
In with the lexicographical order, this set has upper bounds, e.g. It has no least upper bound.
| Mathematics | Order theory | null |
663355 | https://en.wikipedia.org/wiki/Lava%20tube | Lava tube | A lava tube, or pyroduct, is a natural conduit formed by flowing lava from a volcanic vent that moves beneath the hardened surface of a lava flow. If lava in the tube empties, it will leave a cave.
Scientists have discovered fossils in lava tubes near Medina in northwestern Saudi Arabia, encouraging the belief that both humans and animals used the site for refuge from the harsh climate.
Formation
A lava tube is a type of lava cave formed when a low-viscosity lava flow develops a continuous and hard crust, which thickens and forms a roof above the still-flowing lava stream. Tubes form in one of two ways: either by the crusting over of lava channels, or from pāhoehoe flows where the lava is moving under the surface.
Lava usually leaves the point of eruption in channels. These channels tend to stay very hot as their surroundings cool. This means they slowly develop walls around them as the surrounding lava cools and/or as the channel melts its way deeper. These channels can get deep enough to crust over, forming an insulating tube that keeps the lava molten and serves as a conduit for the flowing lava. These types of lava tubes tend to be closer to the lava eruption point.
Farther away from the eruption point, lava can flow in an unchanneled, fan-like manner as it leaves its source, which is usually another lava tube leading back to the eruption point. Called pāhoehoe flows, these areas of surface-moving lava cool, forming either a smooth or rough, ropy surface. The lava continues to flow this way until it begins to block its source. At this point, the subsurface lava is still hot enough to break out at a point, and from this point the lava begins as a new "source". Lava flows from the previous source to this breakout point as the surrounding lava of the pāhoehoe flow cools. This forms an underground channel that becomes a lava tube.
Characteristics
A broad lava-flow field often consists of a main lava tube and a series of smaller tubes that supply lava to the front of one or more separate flows. When the supply of lava stops at the end of an eruption or lava is diverted elsewhere, lava in the tube system drains downslope and leaves partially empty caves.
Such drained tubes commonly exhibit step marks on their walls that mark the various depths at which the lava flowed, known as flow ledges or flow lines depending on how prominently they protrude from the walls. Lava tubes generally have pāhoehoe floors, although this may often be covered in breakdown from the ceiling. A variety of speleothems may be found in lava tubes including a variety of stalactite forms generally known as lavacicles, which can be of the splash, "shark tooth", or tubular varieties. Lavacicles are the most common of lava tube speleothems. Drip stalagmites may form under tubular lava stalactites, and the latter may grade into a form known as a tubular lava helictite. A runner is a bead of lava that is extruded from a small opening and then runs down a wall. Lava tubes may also contain mineral deposits that most commonly take the form of crusts or small crystals, and less commonly, as stalactites and stalagmites. Some stalagmites may contain a central conduit and are interpreted as hornitos extruded from the tube floor.
Lava tubes can be up to wide, though are often narrower, and run anywhere from below the surface. Lava tubes can also be extremely long; one tube from the Mauna Loa 1859 flow enters the ocean about from its eruption point, and the Cueva del Viento–Sobrado system on Teide, Tenerife island, is over long, due to extensive braided maze areas at the upper zones of the system.
A lava tube system in Kiama, Australia, consists of over 20 tubes, many of which are breakouts of a main lava tube. The largest of these lava tubes is in diameter and has columnar jointing due to the large cooling surface. Other tubes have concentric and radial jointing features. The tubes are infilled due to the low slope angle of emplacement.
Extraterrestrial lava tubes
Lunar lava tubes have been discovered and have been studied as possible human habitats, providing natural shielding from radiation. Several holes on the lunar surface, including one in the Marius Hills region, have been observed with angled satellite imagery to lead into voids wider than the holes themselves. These are considered as possible collapses into lunar lava tubes.
Martian lava tubes are associated with innumerable lava flows and lava channels on the flanks of Olympus Mons. Partially collapsed lava tubes are visible as chains of pit craters, and broad lava fans formed by lava emerging from intact, subsurface tubes are also common. Evidence of Martian lava tubes has also been observed on the Southeast Tharsis region and Alba Mons.
Caves, including lava tubes, are considered candidate biotopes of interest for extraterrestrial life.
Notable examples
Iceland
Surtshellir – For a long time, this was the longest known lava tube in the world.
Kenya
Leviathan Cave – At 12.5 kilometres, it is the longest lava tube in Africa.
South Korea
Manjang Cave – more than 8 kilometers long, located in Jeju Island, is a popular tourism spot.
United States
Kazumura Cave, Hawaii – Not only the world's most extensive lava tube, but at , it has the greatest linear extent of any cave known.
| Physical sciences | Caves | Earth science |
663464 | https://en.wikipedia.org/wiki/Buttress | Buttress | A buttress is an architectural structure built against or projecting from a wall which serves to support or reinforce the wall. Buttresses are fairly common on more ancient (typically Gothic) buildings, as a means of providing support to act against the lateral (sideways) forces arising out of inadequately braced roof structures.
The term counterfort can be synonymous with buttress and is often used when referring to dams, retaining walls and other structures holding back earth.
Early examples of buttresses are found on the Eanna Temple (ancient Uruk), dating to as early as the 4th millennium BC.
Terminology
In addition to flying and ordinary buttresses, brick and masonry buttresses that support wall corners can be classified according to their ground plan. A clasping or clamped buttress has an L-shaped ground plan surrounding the corner, an angled buttress has two buttresses meeting at the corner, a setback buttress is similar to an angled buttress but the buttresses are set back from the corner, and a diagonal (or 'French') buttress bisects the angle between the walls where they meet.
The gallery below shows top-down views of various types of buttress (dark grey) supporting the corner wall of a structure (light grey).
Gallery
| Technology | Architectural elements | null |
663616 | https://en.wikipedia.org/wiki/Lion%27s%20mane%20jellyfish | Lion's mane jellyfish | The lion's mane jellyfish (Cyanea capillata) is one of the largest known species of jellyfish. Its range is confined to cold, boreal waters of the Arctic, northern Atlantic, and northern Pacific Oceans. It is common in the English Channel, Irish Sea, North Sea, and in western Scandinavian waters south to Kattegat and Øresund. It may also drift into the southwestern part of the Baltic Sea (where it cannot breed due to the low salinity). Similar jellyfish – which may be the same species – are known to inhabit seas near Australia and New Zealand. The largest recorded specimen was measured off the coast of Massachusetts in 1865 and had a bell with a diameter of and tentacles around long. Lion's mane jellyfish have been observed below 42°N latitude for some time in the larger bays of the East Coast of the United States.
Name
The lion's mane jellyfish is also known as the arctic red jellyfish, hair jelly, snottie, sea blubber or giant jellyfish.
Taxonomy
The taxonomy of the Cyanea species is not fully agreed upon; some zoologists have suggested that all species within the genus should be treated as one. Two distinct taxa, however, occur together in at least the eastern North Atlantic, with the blue jellyfish (Cyanea lamarckii Péron & Lesueur, 1810) differing in color (blue, not red) and smaller size ( diameter, rarely ). Populations in the western Pacific around Japan are sometimes distinguished as Cyanea nozakii, or as a subspecies, C. c. nozakii. In 2015, Russian researchers announced a possible sister species, Cyanea tzetlinii found in the White Sea.
Description
Lion's mane jellyfish (Cyanea capillata) are named for their showy, trailing tentacles reminiscent of a lion's mane. They can vary greatly in size: although capable of attaining a bell diameter of over , those found in lower latitudes are normally smaller than their far northern counterparts, with a bell about in diameter. Furthermore, larger specimens are typically further offshore than smaller ones. Juveniles are lighter orange or tan, very young lion's manes are occasionally colorless and adults are red and start to darken as they age.
While most jellyfish have a circular bell, the bell of the Lion's Mane is divided into eight lobes, giving it the look of an eight-point star. Each lobe contains about 70 to 150 tentacles, arranged in four fairly distinct rows. Along the bell margin is a balance organ at each of the eight indentations between the lobes – the rhopalium – which helps the jellyfish orient itself. From the central mouth extend broad frilly oral arms with many stinging cells called Cnidocytes. Closer to its mouth, its total number of tentacles is around 1,200.
The long, thin tentacles which emanate from the bell's subumbrella have been characterised as "extremely sticky"; they also have stinging cells. The tentacles of larger specimens may trail as long as or more, with the tentacles of the longest known specimen measured at in length, although it has been suggested that this specimen may actually have belonged to a different Cyanea species. This unusual length – longer than a blue whale – has earned it the status of one of the longest known animals in the world.
Behavior and reproduction
Lion's mane jellyfish remain mostly very near the surface, at no more than depth. Their slow pulsations weakly drive them forward, so they depend on ocean currents to travel great distances. The jellyfish are most often spotted during the late summer and autumn, when they have grown to a large size and the currents begin to sweep them to shore. Unlike most pelagic jellyfish, they are completely solitary and rarely travel in groups.
The lion's mane jellyfish uses its stinging tentacles to capture, pull in, and eat prey such as fish, zooplankton, sea creatures, and smaller jellyfish.
Like other jellyfish, lion's manes are capable of both sexual reproduction in the medusa stage and asexual reproduction in the polyp stage. Lion's mane jellyfish have four different stages in their year-long lifespan: a larval stage, a polyp stage, an ephyrae stage, and the medusa stage. The female jellyfish carries its fertilized eggs in a tentacle, where the eggs grow into larvae. When the larvae are old enough, the female deposits them on a hard surface, where the larvae soon grow into polyps. The polyps begin to reproduce asexually, creating stacks of small, immature medusae called ephyrae. The individual ephyrae break off from the stacks, where they eventually grow into the mature medusa stage to become full-grown jellyfish.
Sting and human contact
Human encounters with the jellyfish can cause temporary pain and localized redness. In normal circumstances, however, and in healthy individuals, the stings of the jellyfish are not known to be fatal; vinegar can be used to deactivate the nematocysts. If there is contact with a large number of tentacles, however, medical attention is recommended after exposure.
There may be a significant difference between touching a few tentacles with fingertips at a beach and accidentally swimming into the jellyfish. The initial sensation is more strange than painful and feels like swimming into warmer and somewhat effervescent water. Some minor pain will soon follow. Normally, there is no real danger to humans (with the exception of people suffering from special allergies), but in cases when someone has been stung over large parts of their body by not just the longest tentacles but the entire jellyfish (including the inner tentacles, of which there are around 1,200),
medical attention is recommended as systemic effects can be present. Although rare, severe stings in deep water can also cause panic followed by drowning.
In July 2010, around 150 beachgoers were stung by the remains of a broken-up lion's mane jellyfish along Wallis Sands State Beach in Rye, New Hampshire, US. Considering the size of the species, it is believed that this incident was caused by a single specimen.
In popular culture
A photograph widely distributed on the internet appears to show an anomalously large lion's mane dwarfing a nearby diver. The photo was subsequently shown to be a hoax.
On the television program QI, the show claimed that the longest animal in the world was the lion's mane jellyfish. This was later corrected – in 1864, a bootlace worm (Lineus longissimus) was found washed up on the coast of Fife, Great Britain, that was long. However, this claim is disputed, as bootlace worms can easily stretch to several times their natural length, so it is possible the worm did not actually grow to be that length. If that is the case, the lion's mane jellyfish would indeed be the longest animal in the world.
The creature is referred to as the culprit in the 1926 short story "The Adventure of the Lion's Mane" in The Case-Book of Sherlock Holmes by Sir Arthur Conan Doyle.
Predators
Seabirds, larger fish such as ocean sunfish, other jellyfish species, and most sea turtles will only attack juveniles or smaller specimens while a fully grown adult is incapable of being eaten, due to their massive size and the abundance of stinging tentacles they possess, although both adults and juveniles have been documented eaten by anemones. The leatherback sea turtle feeds almost exclusively on them in large quantities during the summer season around Eastern Canada.
Gallery
| Biology and health sciences | Cnidarians | Animals |
663861 | https://en.wikipedia.org/wiki/Private%20network | Private network | In Internet networking, a private network is a computer network that uses a private address space of IP addresses. These addresses are commonly used for local area networks (LANs) in residential, office, and enterprise environments. Both the IPv4 and the IPv6 specifications define private IP address ranges.
Most Internet service providers (ISPs) allocate only a single publicly routable IPv4 address to each residential customer, but many homes have more than one computer, smartphone, or other Internet-connected device. In this situation, a network address translator (NAT/PAT) gateway is usually used to provide Internet connectivity to multiple hosts. Private addresses are also commonly used in corporate networks which, for security reasons, are not connected directly to the Internet. Often a proxy, SOCKS gateway, or similar devices are used to provide restricted Internet access to network-internal users.
Private network addresses are not allocated to any specific organization. Anyone may use these addresses without approval from regional or local Internet registries. Private IP address spaces were originally defined to assist in delaying IPv4 address exhaustion. IP packets originating from or addressed to a private IP address cannot be routed through the public Internet.
Private addresses are often seen as enhancing network security for the internal network since use of private addresses internally makes it difficult for an external host to initiate a connection to an internal system.
Private IPv4 addresses
The Internet Engineering Task Force (IETF) has directed the Internet Assigned Numbers Authority (IANA) to reserve the following IPv4 address ranges for private networks:
In practice, it is common to subdivide these ranges into smaller subnets.
Dedicated space for carrier-grade NAT deployment
In April 2012, IANA allocated the block of IPv4 addresses specifically for use in carrier-grade NAT scenarios.
This address block should not be used on private networks or on the public Internet. The size of the address block was selected to be large enough to uniquely number all customer access devices for all of a single operator's points of presence in a large metropolitan area such as Tokyo.
Private IPv6 addresses
The concept of private networks has been extended in the next generation of the Internet Protocol, IPv6, and special address blocks are reserved.
The address block is reserved by IANA for unique local addresses (ULAs). They are unicast addresses, but contain a 40-bit random number in the routing prefix to prevent collisions when two private networks are interconnected. Despite being inherently local in usage, the IPv6 address scope of unique local addresses is global.
The first block defined is , designed for routing blocks, in which users can create multiple subnets, as needed.
Examples:
A former standard proposed the use of site-local addresses in the block, but because of scalability concerns and poor definition of what constitutes a site, its use has been deprecated since September 2004.
Link-local addresses
Another type of private networking uses the link-local address range. The validity of link-local addresses is limited to a single link; e.g. to all computers connected to a switch, or to one wireless network. Hosts on different sides of a network bridge are also on the same link, whereas hosts on different sides of a network router are on different links.
IPv4
In IPv4, the utility of link-local addresses is in zero-configuration networking when Dynamic Host Configuration Protocol (DHCP) services are not available and manual configuration by a network administrator is not desirable. The block was allocated for this purpose. If a host on an IEEE 802 (Ethernet) network cannot obtain a network address via DHCP, an address from to may be assigned pseudorandomly. The standard prescribes that address collisions must be handled gracefully.
IPv6
In IPv6, the block is reserved for IP address autoconfiguration.
The implementation of these link-local addresses is mandatory, as various functions of the IPv6 protocol depend on them.
Loopback interface
A special case of private link-local addresses is the loopback interface. These addresses are private and link-local by definition since packets never leave the host device.
IPv4 reserves the entire class A address block for use as private loopback addresses. IPv6 reserves the single address .
Some are advocating reducing to .
Misrouting
It is common for packets originating in private address spaces to be misrouted onto the Internet. Private networks often do not properly configure DNS services for addresses used internally and attempt reverse DNS lookups for these addresses, causing extra traffic to the Internet root nameservers. The AS112 project attempted to mitigate this load by providing special black hole anycast nameservers for private address ranges which only return negative result codes (not found) for these queries.
Organizational edge routers are usually configured to drop ingress IP traffic for these networks, which can occur either by misconfiguration or from malicious traffic using a spoofed source address. Less commonly, ISP edge routers drop such egress traffic from customers, which reduces the impact to the Internet of such misconfigured or malicious hosts on the customer's network.
Merging private networks
Since the private IPv4 address space is relatively small, many private IPv4 networks unavoidably use the same address ranges. This can create a problem when merging such networks, as some addresses may be duplicated for multiple devices. In this case, networks or hosts must be renumbered, often a time-consuming task or a network address translator must be placed between the networks to translate or masquerade one of the address ranges.
IPv6 defines unique local addresses, providing a very large private address space from which each organization can randomly or pseudo-randomly allocate a 40-bit prefix, each of which allows 65536 organizational subnets. With space for about one trillion (1012) prefixes, it is unlikely that two network prefixes in use by different organizations would be the same, provided each of them was selected randomly, as specified in the standard. When two such private IPv6 networks are connected or merged, the risk of an address conflict is therefore virtually absent.
RFC documents
– Address Allocation for Private Internets
– Observations on the use of Components of the Class A Address Space within the Internet
– The Internet Number Registry System
– IPv4 Address Behaviour Today
– IP Network Address Translator (NAT) Terminology and Considerations
– Traditional IP Network Address Translator (Traditional NAT)
– Special-Use IPv4 Addresses (superseded)
– Deprecating Site Local Addresses
– Dynamic Configuration of IPv4 Link-Local Addresses
– Unique Local IPv6 Unicast Addresses
– Special-Use IPv4 Addresses (superseded)
– Reserved IPv4 Prefix for Shared Address Space
– Special-Purpose IP Address Registries
| Technology | Networks | null |
664488 | https://en.wikipedia.org/wiki/Reciprocal%20lattice | Reciprocal lattice | The reciprocal lattice is a term associated with solids with translational symmetry, and plays a major role in many areas such as X-ray and electron diffraction as well as the energies of electrons in a solid. It emerges from the Fourier transform of the lattice associated with the arrangement of the atoms. The direct lattice or real lattice is a periodic function in physical space, such as a crystal system (usually a Bravais lattice). The reciprocal lattice exists in the mathematical space of spatial frequencies, known as reciprocal space or k space, which is the dual of physical space considered as a vector space, and the reciprocal lattice is the sublattice of that space that is dual to the direct lattice.
In quantum physics, reciprocal space is closely related to momentum space according to the proportionality , where is the momentum vector and is the reduced Planck constant. The reciprocal lattice of a reciprocal lattice is equivalent to the original direct lattice, because the defining equations are symmetrical with respect to the vectors in real and reciprocal space. Mathematically, direct and reciprocal lattice vectors represent covariant and contravariant vectors, respectively.
The reciprocal lattice is the set of all vectors , that are wavevectors of plane waves in the Fourier series of a spatial function whose periodicity is the same as that of a direct lattice . Each plane wave in this Fourier series has the same phase or phases that are differed by multiples of at each direct lattice point (so essentially same phase at all the direct lattice points).
The Brillouin zone is a Wigner–Seitz cell of the reciprocal lattice.
Wave-based description
Reciprocal space
Reciprocal space (also called -space) provides a way to visualize the results of the Fourier transform of a spatial function. It is similar in role to the frequency domain arising from the Fourier transform of a time dependent function; reciprocal space is a space over which the Fourier transform of a spatial function is represented at spatial frequencies or wavevectors of plane waves of the Fourier transform. The domain of the spatial function itself is often referred to as real space. In physical applications, such as crystallography, both real and reciprocal space will often each be two or three dimensional. Whereas the number of spatial dimensions of these two associated spaces will be the same, the spaces will differ in their quantity dimension, so that when the real space has the dimension length (L), its reciprocal space will have inverse length, so L−1 (the reciprocal of length).
Reciprocal space comes into play regarding waves, both classical and quantum mechanical. Because a sinusoidal plane wave with unit amplitude can be written as an oscillatory term , with initial phase , angular wavenumber and angular frequency , it can be regarded as a function of both and (and the time-varying part as a function of both and ). This complementary role of and leads to their visualization within complementary spaces (the real space and the reciprocal space). The spatial periodicity of this wave is defined by its wavelength , where ; hence the corresponding wavenumber in reciprocal space will be .
In three dimensions, the corresponding plane wave term becomes , which simplifies to at a fixed time , where is the position vector of a point in real space and now is the wavevector in the three dimensional reciprocal space. (The magnitude of a wavevector is called wavenumber.) The constant is the phase of the wavefront (a plane of a constant phase) through the origin at time , and is a unit vector perpendicular to this wavefront. The wavefronts with phases , where represents any integer, comprise a set of parallel planes, equally spaced by the wavelength .
Reciprocal lattice
In general, a geometric lattice is an infinite, regular array of vertices (points) in space, which can be modelled vectorially as a Bravais lattice. Some lattices may be skew, which means that their primary lines may not necessarily be at right angles. In reciprocal space, a reciprocal lattice is defined as the set of wavevectors of plane waves in the Fourier series of any function whose periodicity is compatible with that of an initial direct lattice in real space. Equivalently, a wavevector is a vertex of the reciprocal lattice if it corresponds to a plane wave in real space whose phase at any given time is the same (actually differs by with an integer ) at every direct lattice vertex.
One heuristic approach to constructing the reciprocal lattice in three dimensions is to write the position vector of a vertex of the direct lattice as , where the are integers defining the vertex and the are linearly independent primitive translation vectors (or shortly called primitive vectors) that are characteristic of the lattice. There is then a unique plane wave (up to a factor of negative one), whose wavefront through the origin contains the direct lattice points at and , and with its adjacent wavefront (whose phase differs by or from the former wavefront passing the origin) passing through . Its angular wavevector takes the form , where is the unit vector perpendicular to these two adjacent wavefronts and the wavelength must satisfy , means that is equal to the distance between the two wavefronts. Hence by construction and .
Cycling through the indices in turn, the same method yields three wavevectors with , where the Kronecker delta equals one when and is zero otherwise. The comprise a set of three primitive wavevectors or three primitive translation vectors for the reciprocal lattice, each of whose vertices takes the form , where the are integers. The reciprocal lattice is also a Bravais lattice as it is formed by integer combinations of the primitive vectors, that are , , and in this case. Simple algebra then shows that, for any plane wave with a wavevector on the reciprocal lattice, the total phase shift between the origin and any point on the direct lattice is a multiple of (that can be possibly zero if the multiplier is zero), so the phase of the plane wave with will essentially be equal for every direct lattice vertex, in conformity with the reciprocal lattice definition above. (Although any wavevector on the reciprocal lattice does always take this form, this derivation is motivational, rather than rigorous, because it has omitted the proof that no other possibilities exist.)
The Brillouin zone is a primitive cell (more specifically a Wigner–Seitz cell) of the reciprocal lattice, which plays an important role in solid state physics due to Bloch's theorem. In pure mathematics, the dual space of linear forms and the dual lattice provide more abstract generalizations of reciprocal space and the reciprocal lattice.
Mathematical description
Assuming a three-dimensional Bravais lattice and labelling each lattice vector (a vector indicating a lattice point) by the subscript as 3-tuple of integers,
where
where is the set of integers and is a primitive translation vector or shortly primitive vector. Taking a function where is a position vector from the origin to any position, if follows the periodicity of this lattice, e.g. the function describing the electronic density in an atomic crystal, it is useful to write as a multi-dimensional Fourier series
where now the subscript , so this is a triple sum.
As follows the periodicity of the lattice, translating by any lattice vector we get the same value, hence
Expressing the above instead in terms of their Fourier series we have
Because equality of two Fourier series implies equality of their coefficients, , which only holds when
where
Mathematically, the reciprocal lattice is the set of all vectors , that are wavevectors of plane waves in the Fourier series of a spatial function whose periodicity is the same as that of a direct lattice as the set of all direct lattice point position vectors , and satisfy this equality for all . Each plane wave in the Fourier series has the same phase (actually can be differed by a multiple of ) at all the lattice point .
As shown in the section multi-dimensional Fourier series, can be chosen in the form of where . With this form, the reciprocal lattice as the set of all wavevectors for the Fourier series of a spatial function which periodicity follows , is itself a Bravais lattice as it is formed by integer combinations of its own primitive translation vectors , and the reciprocal of the reciprocal lattice is the original lattice, which reveals the Pontryagin duality of their respective vector spaces. (There may be other form of . Any valid form of results in the same reciprocal lattice.)
Two dimensions
For an infinite two-dimensional lattice, defined by its primitive vectors , its reciprocal lattice can be determined by generating its two reciprocal primitive vectors, through the following formulae,
where is an integer and
Here represents a 90 degree rotation matrix, i.e. a quarter turn. The anti-clockwise rotation and the clockwise rotation can both be used to determine the reciprocal lattice: If is the anti-clockwise rotation and is the clockwise rotation, for all vectors . Thus, using the permutation
we obtain
Notably, in a 3D space this 2D reciprocal lattice is an infinitely extended set of Bragg rods—described by Sung et al.
Three dimensions
For an infinite three-dimensional lattice , defined by its primitive vectors and the subscript of integers , its reciprocal lattice with the integer subscript can be determined by generating its three reciprocal primitive vectors
where is the scalar triple product. The choice of these is to satisfy as the known condition (There may be other condition.) of primitive translation vectors for the reciprocal lattice derived in the heuristic approach above and the section multi-dimensional Fourier series. This choice also satisfies the requirement of the reciprocal lattice mathematically derived above. Using column vector representation of (reciprocal) primitive vectors, the formulae above can be rewritten using matrix inversion:
This method appeals to the definition, and allows generalization to arbitrary dimensions. The cross product formula dominates introductory materials on crystallography.
The above definition is called the "physics" definition, as the factor of comes naturally from the study of periodic structures. An essentially equivalent definition, the "crystallographer's" definition, comes from defining the reciprocal lattice . which changes the reciprocal primitive vectors to be
and so on for the other primitive vectors. The crystallographer's definition has the advantage that the definition of is just the reciprocal magnitude of in the direction of , dropping the factor of . This can simplify certain mathematical manipulations, and expresses reciprocal lattice dimensions in units of spatial frequency. It is a matter of taste which definition of the lattice is used, as long as the two are not mixed.
is conventionally written as or , called Miller indices; is replaced with , replaced with , and replaced with . Each lattice point in the reciprocal lattice corresponds to a set of lattice planes in the real space lattice. (A lattice plane is a plane crossing lattice points.) The direction of the reciprocal lattice vector corresponds to the normal to the real space planes. The magnitude of the reciprocal lattice vector is given in reciprocal length and is equal to the reciprocal of the interplanar spacing of the real space planes.
Higher dimensions
The formula for dimensions can be derived assuming an -dimensional real vector space with a basis and an inner product . The reciprocal lattice vectors are uniquely determined by the formula . Using the permutation
they can be determined with the following formula:
Here, is the volume form, is the inverse of the vector space isomorphism defined by and denotes the inner multiplication.
One can verify that this formula is equivalent to the known formulas for the two- and three-dimensional case by using the following facts: In three dimensions, and in two dimensions, , where is the rotation by 90 degrees (just like the volume form, the angle assigned to a rotation depends on the choice of orientation).
Reciprocal lattices of various crystals
Reciprocal lattices for the cubic crystal system are as follows.
Simple cubic lattice
The simple cubic Bravais lattice, with cubic primitive cell of side , has for its reciprocal a simple cubic lattice with a cubic primitive cell of side (or in the crystallographer's definition). The cubic lattice is therefore said to be self-dual, having the same symmetry in reciprocal space as in real space.
Face-centered cubic (FCC) lattice
The reciprocal lattice to an FCC lattice is the body-centered cubic (BCC) lattice, with a cube side of .
Consider an FCC compound unit cell. Locate a primitive unit cell of the FCC; i.e., a unit cell with one lattice point. Now take one of the vertices of the primitive unit cell as the origin. Give the basis vectors of the real lattice. Then from the known formulae, you can calculate the basis vectors of the reciprocal lattice. These reciprocal lattice vectors of the FCC represent the basis vectors of a BCC real lattice. The basis vectors of a real BCC lattice and the reciprocal lattice of an FCC resemble each other in direction but not in magnitude.
Body-centered cubic (BCC) lattice
The reciprocal lattice to a BCC lattice is the FCC lattice, with a cube side of .
It can be proven that only the Bravais lattices which have 90 degrees between (cubic, tetragonal, orthorhombic) have primitive translation vectors for the reciprocal lattice, , parallel to their real-space vectors.
Simple hexagonal lattice
The reciprocal to a simple hexagonal Bravais lattice with lattice constants and is another simple hexagonal lattice with lattice constants and rotated through 90° about the c axis with respect to the direct lattice. The simple hexagonal lattice is therefore said to be self-dual, having the same symmetry in reciprocal space as in real space. Primitive translation vectors for this simple hexagonal Bravais lattice vectors are
Arbitrary collection of atoms
One path to the reciprocal lattice of an arbitrary collection of atoms comes from the idea of scattered waves in the Fraunhofer (long-distance or lens back-focal-plane) limit as a Huygens-style sum of amplitudes from all points of scattering (in this case from each individual atom). This sum is denoted by the complex amplitude in the equation below, because it is also the Fourier transform (as a function of spatial frequency or reciprocal distance) of an effective scattering potential in direct space:
Here g = q/(2) is the scattering vector q in crystallographer units, N is the number of atoms, fj[g] is the atomic scattering factor for atom j and scattering vector g, while rj is the vector position of atom j. The Fourier phase depends on one's choice of coordinate origin.
For the special case of an infinite periodic crystal, the scattered amplitude F = M Fh,k,ℓ from M unit cells (as in the cases above) turns out to be non-zero only for integer values of , where
when there are j = 1,m atoms inside the unit cell whose fractional lattice indices are respectively {uj, vj, wj}. To consider effects due to finite crystal size, of course, a shape convolution for each point or the equation above for a finite lattice must be used instead.
Whether the array of atoms is finite or infinite, one can also imagine an "intensity reciprocal lattice" I[g], which relates to the amplitude lattice F via the usual relation I = F*F where F* is the complex conjugate of F. Since Fourier transformation is reversible, of course, this act of conversion to intensity tosses out "all except 2nd moment" (i.e. the phase) information. For the case of an arbitrary collection of atoms, the intensity reciprocal lattice is therefore:
Here rjk is the vector separation between atom j and atom k. One can also use this to predict the effect of nano-crystallite shape, and subtle changes in beam orientation, on detected diffraction peaks even if in some directions the cluster is only one atom thick. On the down side, scattering calculations using the reciprocal lattice basically consider an incident plane wave. Thus after a first look at reciprocal lattice (kinematic scattering) effects, beam broadening and multiple scattering (i.e. dynamical) effects may be important to consider as well.
Generalization of a dual lattice
There are actually two versions in mathematics of the abstract dual lattice concept, for a given lattice L in a real vector space V, of finite dimension.
The first, which generalises directly the reciprocal lattice construction, uses Fourier analysis. It may be stated simply in terms of Pontryagin duality. The dual group V^ to V is again a real vector space, and its closed subgroup L^ dual to L turns out to be a lattice in V^. Therefore, L^ is the natural candidate for dual lattice, in a different vector space (of the same dimension).
The other aspect is seen in the presence of a quadratic form Q on V; if it is non-degenerate it allows an identification of the dual space V* of V with V. The relation of V* to V is not intrinsic; it depends on a choice of Haar measure (volume element) on V. But given an identification of the two, which is in any case well-defined up to a scalar, the presence of Q allows one to speak to the dual lattice to L while staying within V.
In mathematics, the dual lattice of a given lattice L in an abelian locally compact topological group G is the subgroup L∗ of the dual group of G consisting of all continuous characters that are equal to one at each point of L.
In discrete mathematics, a lattice is a locally discrete set of points described by all integral linear combinations of linearly independent vectors in Rn. The dual lattice is then defined by all points in the linear span of the original lattice (typically all of Rn) with the property that an integer results from the inner product with all elements of the original lattice. It follows that the dual of the dual lattice is the original lattice.
Furthermore, if we allow the matrix B to have columns as the linearly independent vectors that describe the lattice, then the matrix has columns of vectors that describe the dual lattice.
| Physical sciences | Crystallography | Physics |
664497 | https://en.wikipedia.org/wiki/Parallel%20%28geometry%29 | Parallel (geometry) | In geometry, parallel lines are coplanar infinite straight lines that do not intersect at any point. Parallel planes are planes in the same three-dimensional space that never meet. Parallel curves are curves that do not touch each other or intersect and keep a fixed minimum distance. In three-dimensional Euclidean space, a line and a plane that do not share a point are also said to be parallel. However, two noncoplanar lines are called skew lines. Line segments and Euclidean vectors are parallel if they have the same direction or opposite direction (not necessarily the same length).
Parallel lines are the subject of Euclid's parallel postulate. Parallelism is primarily a property of affine geometries and Euclidean geometry is a special instance of this type of geometry.
In some other geometries, such as hyperbolic geometry, lines can have analogous properties that are referred to as parallelism.
Symbol
The parallel symbol is . For example, indicates that line AB is parallel to line CD.
In the Unicode character set, the "parallel" and "not parallel" signs have codepoints U+2225 (∥) and U+2226 (∦), respectively. In addition, U+22D5 (⋕) represents the relation "equal and parallel to".
Euclidean parallelism
Two lines in a plane
Conditions for parallelism
Given parallel straight lines l and m in Euclidean space, the following properties are equivalent:
Every point on line m is located at exactly the same (minimum) distance from line l (equidistant lines).
Line m is in the same plane as line l but does not intersect l (recall that lines extend to infinity in either direction).
When lines m and l are both intersected by a third straight line (a transversal) in the same plane, the corresponding angles of intersection with the transversal are congruent.
Since these are equivalent properties, any one of them could be taken as the definition of parallel lines in Euclidean space, but the first and third properties involve measurement, and so, are "more complicated" than the second. Thus, the second property is the one usually chosen as the defining property of parallel lines in Euclidean geometry. The other properties are then consequences of Euclid's Parallel Postulate.
History
The definition of parallel lines as a pair of straight lines in a plane which do not meet appears as Definition 23 in Book I of Euclid's Elements. Alternative definitions were discussed by other Greeks, often as part of an attempt to prove the parallel postulate. Proclus attributes a definition of parallel lines as equidistant lines to Posidonius and quotes Geminus in a similar vein. Simplicius also mentions Posidonius' definition as well as its modification by the philosopher Aganis.
At the end of the nineteenth century, in England, Euclid's Elements was still the standard textbook in secondary schools. The traditional treatment of geometry was being pressured to change by the new developments in projective geometry and non-Euclidean geometry, so several new textbooks for the teaching of geometry were written at this time. A major difference between these reform texts, both between themselves and between them and Euclid, is the treatment of parallel lines. These reform texts were not without their critics and one of them, Charles Dodgson (a.k.a. Lewis Carroll), wrote a play, Euclid and His Modern Rivals, in which these texts are lambasted.
One of the early reform textbooks was James Maurice Wilson's Elementary Geometry of 1868. Wilson based his definition of parallel lines on the primitive notion of direction. According to Wilhelm Killing the idea may be traced back to Leibniz. Wilson, without defining direction since it is a primitive, uses the term in other definitions such as his sixth definition, "Two straight lines that meet one another have different directions, and the difference of their directions is the angle between them." In definition 15 he introduces parallel lines in this way; "Straight lines which have the same direction, but are not parts of the same straight line, are called parallel lines." Augustus De Morgan reviewed this text and declared it a failure, primarily on the basis of this definition and the way Wilson used it to prove things about parallel lines. Dodgson also devotes a large section of his play (Act II, Scene VI § 1) to denouncing Wilson's treatment of parallels. Wilson edited this concept out of the third and higher editions of his text.
Other properties, proposed by other reformers, used as replacements for the definition of parallel lines, did not fare much better. The main difficulty, as pointed out by Dodgson, was that to use them in this way required additional axioms to be added to the system. The equidistant line definition of Posidonius, expounded by Francis Cuthbertson in his 1874 text Euclidean Geometry suffers from the problem that the points that are found at a fixed given distance on one side of a straight line must be shown to form a straight line. This can not be proved and must be assumed to be true. The corresponding angles formed by a transversal property, used by W. D. Cooley in his 1860 text, The Elements of Geometry, simplified and explained requires a proof of the fact that if one transversal meets a pair of lines in congruent corresponding angles then all transversals must do so. Again, a new axiom is needed to justify this statement.
Construction
The three properties above lead to three different methods of construction of parallel lines.
Distance between two parallel lines
Because parallel lines in a Euclidean plane are equidistant there is a unique distance between the two parallel lines. Given the equations of two non-vertical, non-horizontal parallel lines,
the distance between the two lines can be found by locating two points (one on each line) that lie on a common perpendicular to the parallel lines and calculating the distance between them. Since the lines have slope m, a common perpendicular would have slope −1/m and we can take the line with equation y = −x/m as a common perpendicular. Solve the linear systems
and
to get the coordinates of the points. The solutions to the linear systems are the points
and
These formulas still give the correct point coordinates even if the parallel lines are horizontal (i.e., m = 0). The distance between the points is
which reduces to
When the lines are given by the general form of the equation of a line (horizontal and vertical lines are included):
their distance can be expressed as
Two lines in three-dimensional space
Two lines in the same three-dimensional space that do not intersect need not be parallel. Only if they are in a common plane are they called parallel; otherwise they are called skew lines.
Two distinct lines l and m in three-dimensional space are parallel if and only if the distance from a point P on line m to the nearest point on line l is independent of the location of P on line m. This never holds for skew lines.
A line and a plane
A line m and a plane q in three-dimensional space, the line not lying in that plane, are parallel if and only if they do not intersect.
Equivalently, they are parallel if and only if the distance from a point P on line m to the nearest point in plane q is independent of the location of P on line m.
Two planes
Similar to the fact that parallel lines must be located in the same plane, parallel planes must be situated in the same three-dimensional space and contain no point in common.
Two distinct planes q and r are parallel if and only if the distance from a point P in plane q to the nearest point in plane r is independent of the location of P in plane q. This will never hold if the two planes are not in the same three-dimensional space.
In non-Euclidean geometry
In non-Euclidean geometry, the concept of a straight line is replaced by the more general concept of a geodesic, a curve which is locally straight with respect to the metric (definition of distance) on a Riemannian manifold, a surface (or higher-dimensional space) which may itself be curved. In general relativity, particles not under the influence of external forces follow geodesics in spacetime, a four-dimensional manifold with 3 spatial dimensions and 1 time dimension.
In non-Euclidean geometry (elliptic or hyperbolic geometry) the three Euclidean properties mentioned above are not equivalent and only the second one (Line m is in the same plane as line l but does not intersect l) is useful in non-Euclidean geometries, since it involves no measurements. In general geometry the three properties above give three different types of curves, equidistant curves, parallel geodesics and geodesics sharing a common perpendicular, respectively.
Hyperbolic geometry
While in Euclidean geometry two geodesics can either intersect or be parallel, in hyperbolic geometry, there are three possibilities. Two geodesics belonging to the same plane can either be:
intersecting, if they intersect in a common point in the plane,
parallel, if they do not intersect in the plane, but converge to a common limit point at infinity (ideal point), or
ultra parallel, if they do not have a common limit point at infinity.
In the literature ultra parallel geodesics are often called non-intersecting. Geodesics intersecting at infinity are called limiting parallel.
As in the illustration through a point a not on line l there are two limiting parallel lines, one for each direction ideal point of line l. They separate the lines intersecting line l and those that are ultra parallel to line l.
Ultra parallel lines have single common perpendicular (ultraparallel theorem), and diverge on both sides of this common perpendicular.
Spherical or elliptic geometry
In spherical geometry, all geodesics are great circles. Great circles divide the sphere in two equal hemispheres and all great circles intersect each other. Thus, there are no parallel geodesics to a given geodesic, as all geodesics intersect. Equidistant curves on the sphere are called parallels of latitude analogous to the latitude lines on a globe. Parallels of latitude can be generated by the intersection of the sphere with a plane parallel to a plane through the center of the sphere.
Reflexive variant
If l, m, n are three distinct lines, then
In this case, parallelism is a transitive relation. However, in case l = n, the superimposed lines are not considered parallel in Euclidean geometry. The binary relation between parallel lines is evidently a symmetric relation. According to Euclid's tenets, parallelism is not a reflexive relation and thus fails to be an equivalence relation. Nevertheless, in affine geometry a pencil of parallel lines is taken as an equivalence class in the set of lines where parallelism is an equivalence relation.
To this end, Emil Artin (1957) adopted a definition of parallelism where two lines are parallel if they have all or none of their points in common.
Then a line is parallel to itself so that the reflexive and transitive properties belong to this type of parallelism, creating an equivalence relation on the set of lines. In the study of incidence geometry, this variant of parallelism is used in the affine plane.
| Mathematics | Two-dimensional space | null |
664655 | https://en.wikipedia.org/wiki/Gal%20%28unit%29 | Gal (unit) | The gal (symbol: Gal), sometimes called galileo after Galileo Galilei, is a unit of acceleration typically used in precision gravimetry. The gal is defined as 1 centimeter per second squared (1 cm/s2). The milligal (mGal) and microgal (μGal) are respectively one thousandth and one millionth of a gal.
The gal is not part of the International System of Units (known by its French-language initials "SI"). In 1978 the CIPM decided that it was permissible to use the gal "with the SI until the CIPM considers that [its] use is no longer necessary". Use of the gal was deprecated by the standard ISO 80000-3:2006, now superseded.
The gal is a derived unit, defined in terms of the centimeter–gram–second (CGS) base unit of length, the centimeter, and the second, which is the base unit of time in both the CGS and the modern SI system. In SI base units, 1 Gal is equal to 0.01 m/s2.
The acceleration due to Earth's gravity at its surface is 976 to 983 Gal, the variation being due mainly to differences in latitude and elevation. Standard gravity is 980.665 Gal. Mountains and masses of lesser density within the Earth's crust typically cause variations in gravitational acceleration of tens to hundreds of milligals (mGal). The gravity gradient (variation with height) above Earth's surface is about 3.1 μGal per centimeter of height (), resulting in a maximal difference of about 2 Gal (0.02 m/s2) from the top of Mount Everest to sea level.
Unless it is being used at the beginning of a sentence or in paragraph or section titles, the unit name gal is properly spelled with a lowercase g. As with the torr and its symbol, the unit name (gal) and its symbol (Gal) are spelled identically except that the latter is capitalized.
Conversions
| Physical sciences | Acceleration | Basics and measurement |
12264442 | https://en.wikipedia.org/wiki/Sea%20surface%20microlayer | Sea surface microlayer | The sea surface microlayer (SML) is the boundary interface between the atmosphere and ocean, covering about 70% of Earth's surface. With an operationally defined thickness between 1 and , the SML has physicochemical and biological properties that are measurably distinct from underlying waters. Recent studies now indicate that the SML covers the ocean to a significant extent, and evidence shows that it is an aggregate-enriched biofilm environment with distinct microbial communities. Because of its unique position at the air-sea interface, the SML is central to a range of global marine biogeochemical and climate-related processes.
The sea surface microlayer is the boundary layer where all exchange occurs between the atmosphere and the ocean. The chemical, physical, and biological properties of the SML differ greatly from the sub-surface water just a few centimeters beneath.
Despite the huge extent of the ocean's surface, until now relatively little attention has been paid to the sea surface microlayer (SML) as the ultimate interface where heat, momentum and mass exchange between the ocean and the atmosphere takes place. Via the SML, large-scale environmental changes in the ocean such as warming, acidification, deoxygenation, and eutrophication potentially influence cloud formation, precipitation, and the global radiation balance. Due to the deep connectivity between biological, chemical, and physical processes, studies of the SML may reveal multiple sensitivities to global and regional changes.
Understanding the processes at the ocean's surface, in particular involving the SML as an important and determinant interface, could provide an essential contribution to the reduction of uncertainties regarding ocean-climate feedbacks. As of 2017, processes occurring within the SML, as well as the associated rates of material exchange through the SML, remained poorly understood and were rarely represented in marine and atmospheric numerical models.
Overview
The sea surface microlayer (SML) is the boundary interface between the atmosphere and ocean, covering about 70% of the Earth's surface. The SML has physicochemical and biological properties that are measurably distinct from underlying waters. Because of its unique position at the air-sea interface, the SML is central to a range of global biogeochemical and climate-related processes. Although known for the last six decades, the SML often has remained in a distinct research niche, primarily as it was not thought to exist under typical oceanic conditions. Recent studies now indicate that the SML covers the ocean to a significant extent, highlighting its global relevance as the boundary layer linking two major components of the Earth system – the ocean and the atmosphere.
In 1983, Sieburth hypothesised that the SML was a hydrated gel-like layer formed by a complex mixture of carbohydrates, proteins, and lipids. In recent years, his hypothesis has been confirmed, and scientific evidence indicates that the SML is an aggregate-enriched biofilm environment with distinct microbial communities. In 1999 Ellison et al. estimated that 200 Tg C yr−1 (200 million tonnes of carbon per year) accumulates in the SML, similar to sedimentation rates of carbon to the ocean's seabed, though the accumulated carbon in the SML probably has a very short residence time. Although the total volume of the microlayer is very small compared to the ocean's volume, Carlson suggested in his seminal 1993 paper that unique interfacial reactions may occur in the SML that may not occur in the underlying water or at a much slower rate there. He therefore hypothesised that the SML plays an important role in the diagenesis of carbon in the upper ocean. Biofilm-like properties and highest possible exposure to solar radiation leads to an intuitive assumption that the SML is a biochemical microreactor.
Historically, the SML has been summarized as being a microhabitat composed of several layers distinguished by their ecological, chemical and physical properties with an operational total thickness of between 1 and 1000 μm. In 2005 Hunter defined the SML as a "microscopic portion of the surface ocean which is in contact with the atmosphere and which may have physical, chemical or biological properties that are measurably different from those of adjacent sub-surface waters". He avoids a definite range of thickness as it depends strongly on the feature of interest. A thickness of 60 μm has been measured based on sudden changes of the pH, and could be meaningfully used for studying the physicochemical properties of the SML. At such thickness, the SML represents a laminar layer, free of turbulence, and greatly affecting the exchange of gases between the ocean and atmosphere. As a habitat for neuston (surface-dwelling organisms ranging from bacteria to larger siphonophores), the thickness of the SML in some ways depends on the organism or ecological feature of interest. In 2005, Zaitsev described the SML and associated near-surface layer (down to 5 cm) as an incubator or nursery for eggs and larvae for a wide range of aquatic organisms.
Hunter's definition includes all interlinked layers from the laminar layer to the nursery without explicit reference to defined depths. In 2017, Wurl et al. proposed Hunter's definition be validated with a redeveloped SML paradigm that includes its global presence, biofilm-like properties and role as a nursery. The new paradigm pushes the SML into a new and wider context relevant to many ocean and climate sciences.
According to Wurl et al., the SML can never be devoid of organics due to the abundance of surface-active substances (e.g., surfactants) in the upper ocean and the phenomenon of surface tension at air-liquid interfaces. The SML is analogous to the thermal boundary layer, and remote sensing of the sea surface temperature shows ubiquitous anomalies between the sea surface skin and bulk temperature. Even so, the differences in both are driven by different processes. Enrichment, defined as concentration ratios of an analyte in the SML to the underlying bulk water, has been used for decades as evidence for the existence of the SML. Consequently, depletions of organics in the SML are debatable; however, the question of enrichment or depletion is likely to be a function of the thickness of the SML (which varies with sea state; including losses via sea spray, the concentrations of organics in the bulk water, and the limitations of sampling techniques to collect thin layers . Enrichment of surfactants, and changes in the sea surface temperature and salinity, serve as universal indicators for the presence of the SML. Organisms are perhaps less suitable as indicators of the SML because they can actively avoid the SML and/or the harsh conditions in the SML may reduce their populations. However, the thickness of the SML remains "operational" in field experiments because the thickness of the collected layer is governed by the sampling method. Advances in SML sampling technology are needed to improve our understanding of how the SML influences air-sea interactions.
Marine surface habitats sit at the interface between the atmosphere and the ocean. The biofilm-like habitat at the surface of the ocean harbours surface-dwelling microorganisms, commonly referred to as neuston.
The sea surface microlayer (SML) constitutes the uppermost layer of the ocean, only 1–1000 μm thick, with unique chemical and biological properties that distinguish it from the underlying water (ULW). Due to the location at the air-sea interface, the SML can influence exchange processes across this boundary layer, such as air-sea gas exchange and the formation of sea spray aerosols.
Due to its exclusive position between the atmosphere and the hydrosphere and by spanning about 70% of the Earth's surface, the sea-surface microlayer (sea-SML) is regarded as a fundamental component in air–sea exchange processes and in biogeochemical cycling. Although having a minor thickness of <1000 μm, the elusive SML is long known for its distinct physicochemical characteristics compared to the underlying water, e.g., by featuring the accumulation of dissolved and particulate organic matter, transparent exopolymer particles (TEP), and surface-active molecules. Therefore, the SML is a gelatinous biofilm, maintaining physical stability through surface tension forces. It also forms a vast habitat for different organisms, collectively termed as neuston with a recent global estimate of 2 × 1023 microbial cells for the sea-SML.
Life at air–water interfaces has never been considered easy, mainly because of the harsh environmental conditions that influence the SML. However, high abundances of microorganisms, especially of bacteria and picophytoplankton, accumulating in the SML compared to the underlying water were frequently reported, accompanied by a predominant heterotrophic activity. This is because primary production at the immediate air–water interface is often hindered by photoinhibition. However, some exceptions of photosynthetic organisms, e.g., Trichodesmium, Synechococcus, or Sargassum, show more tolerance towards high light intensities and, hence, can become enriched in the SML. Previous research has provided evidence that neustonic organisms can cope with wind and wave energy, solar and ultraviolet (UV) radiation, fluctuations in temperature and salinity, and a higher potential predation risk by the zooneuston. Furthermore, wind action promoting sea spray formation and bubbles rising from deeper water and bursting at the surface release SML-associated microbes into the atmosphere. In addition to being more concentrated compared to planktonic counterparts, the bacterioneuston, algae, and protists display distinctive community compositions compared to the underlying water, in both marine and freshwater habitats. Furthermore, the bacterial community composition was often dependent on the SML sampling device being used. While being well defined with respect to bacterial community composition, little is known about viruses in the SML, i.e., the virioneuston. This review has its focus on virus–bacterium dynamics at air–water interfaces, even if viruses likely interact with other SML microbes, including archaea and the phytoneuston, as can be deduced from viral interference with their planktonic counterparts. Although viruses were briefly mentioned as pivotal SML components in a recent review on this unique habitat, a synopsis of the emerging knowledge and the major research gaps regarding bacteriophages at air–water interfaces is still missing in the literature.
Properties
Organic compounds such as amino acids, carbohydrates, fatty acids, and phenols are highly enriched in the SML interface. Most of these come from biota in the sub-surface waters, which decay and become transported to the surface, though other sources exist also such as atmospheric deposition, coastal runoff, and anthropogenic nutrification. The relative concentration of these compounds is dependent on the nutrient sources as well as climate conditions such as wind speed and precipitation. These organic compounds on the surface create a "film," referred to as a "slick" when visible, which affects the physical and optical properties of the interface. These films occur because of the hydrophobic tendencies of many organic compounds, which causes them to protrude into the air-interface. The existence of organic surfactants on the ocean surface impedes wave formation for low wind speeds. For increasing concentrations of surfactant there is an increasing critical wind speed necessary to create ocean waves. Increased levels of organic compounds at the surface also hinders air-sea gas exchange at low wind speeds. One way in which particulates and organic compounds on the surface are transported into the atmosphere is the process called "bubble bursting". Bubbles generate the major portion of marine aerosols. They can be dispersed to heights of several meters, picking up whatever particles latch on to their surface. However, the major supplier of materials comes from the SML.
Processes
Surfaces and interfaces are critical zones where major physical, chemical, and biological exchanges occur. As the ocean covers 362 million km2, about 71% of the Earth's surface, the ocean-atmosphere interface is plausibly one of the largest and most important interfaces on the planet. Every substance entering or leaving the ocean from or to the atmosphere passes through this interface, which on the water-side -and to a lesser extent on the air-side- shows distinct physical, chemical, and biological properties. On the water side the uppermost 1 to 1000 μm of this interface are referred to as the sea surface microlayer (SML). Like a skin, the SML is expected to control the rates of exchange of energy and matter between air and sea, thereby potentially exerting both short-term and long-term impacts on various Earth system processes, including biogeochemical cycling, production and uptake of radiately active gases like or DMS, thus ultimately climate regulation. As of 2017, processes occurring within the SML, as well as the associated rates of material exchange through the SML, remained poorly understood and were rarely represented in marine and atmospheric numerical models.
An improved understanding of the biological, chemical, and physical processes at the ocean's upper surface could provide an essential contribution to the reduction of uncertainties regarding ocean-climate feedbacks. Due to its positioning between atmosphere and ocean, the SML is the first to be exposed to climate changes including temperature, climate relevant trace gases, wind speed, and precipitation as well as to pollution by human waste, including nutrients, toxins, nanomaterials, and plastic debris.
Bacterioneuston
The term neuston describes the organisms in the SML and was first suggested by Naumann in 1917. As in other marine ecosystems, bacterioneuston communities have important roles in SML functioning. Bacterioneuston community composition of the SML has been analysed and compared to the underlying water in different habitats with varying results, and has primarily focused on coastal waters and shelf seas, with limited study of the open ocean . In the North Sea, a distinct bacterial community was found in the SML with Vibrio spp. and Pseudoalteromonas spp. dominating the bacterioneuston. During an artificially induced phytoplankton bloom in a fjord mesocosm experiment, the most dominant denaturing gradient gel electrophoresis (DGGE) bands of the bacterioneuston consisted of two bacterial families: Flavobacteriaceae and Alteromonadaceae. Other studies have however, found little or no differences in the bacterial community composition of the SML and the ULW. Difficulties in direct comparisons between studies can arise because of the different methods used to sample the SML, which result in varied sampling depths.
Even less is known about the community control mechanisms in the SML and how the bacterial community assembles at the air-sea interface. The bacterioneuston community could be altered by differing wind conditions and radiation levels, with high wind speeds inhibiting the formation of a distinct bacterioneuston community. Wind speed and radiation levels refer to external controls, however, bacterioneuston community composition might also be influenced by internal factors such as nutrient availability and organic matter (OM) produced either in the SML or in the ULW.
One of the principal OM components consistently enriched in the SML are transparent exopolymer particles (TEP), which are rich in carbohydrates and form by the aggregation of dissolved precursors excreted by phytoplankton in the euphotic zone. Higher TEP formation rates in the SML, facilitated through wind shear and dilation of the surface water, have been proposed as one explanation for the observed enrichment in TEP. Also, due to their natural positive buoyancy, when not ballasted by other particles sticking to them, TEP ascend through the water column and ultimately end up at the SML . A second possible pathway of TEP from the water column to the SML is by bubble scavenging.
Next to rising bubbles, another potential transport mechanism for bacteria from the ULW to the SML could be ascending particles or more specifically TEP. Bacteria readily attach to TEP in the water column. TEP can serve as microbial hotspots and can be used directly as a substrate for bacterial degradation, and as grazing protection for attached bacteria, e.g., by acting as an alternate food source for zooplankton. TEP have also been suggested to serve as light protection for microorganisms in environments with high irradiation.
Virioneuston
Viruses in the sea surface microlayer, the so-called virioneuston, have recently become of interest to researchers as enigmatic biological entities in the boundary surface layers with potentially important ecological impacts. Given this vast air–water interface sits at the intersection of major air–water exchange processes spanning more than 70% of the global surface area, it is likely to have profound implications for marine biogeochemical cycles, on the microbial loop and gas exchange, as well as the marine food web structure, the global dispersal of airborne viruses originating from the sea surface microlayer, and human health.
Viruses are the most abundant biological entities in the water column of the world's oceans. In the free water column, the virioplankton typically outnumbers the bacterioplankton by one order of magnitude reaching typical bulk water concentrations of 107 viruses mL−1. Moreover, they are known as integral parts of global biogeochemical cycles to shape and drive microbial diversity and to structure trophic networks. Like other neuston members, the virioneuston likely originates from the bulk seawater. For instance, in 1977 Baylor et al. postulated adsorption of viruses onto air bubbles as they rise to the surface, or viruses can stick to organic particles also being transported to the SML via bubble scavenging.
Within the SML, viruses interacting with the bacterioneuston will probably induce the viral shunt, a phenomenon that is well known for marine pelagic systems. The term viral shunt describes the release of organic carbon and other nutritious compounds from the virus-mediated lysis of host cells, and its addition to the local dissolved organic matter (DOM) pool. The enriched and densely packed bacterioneuston forms an excellent target for viruses compared to the bacterioplankton populating the subsurface. This is because high host-cell numbers will increase the probability of host–virus encounters. The viral shunt might effectively contribute to the SML's already high DOM content enhancing bacterial production as previously suggested for pelagic ecosystems and in turn replenishing host cells for viral infections. By affecting the DOM pool, viruses in the SML might directly interfere with the microbial loop being initiated when DOM is microbially recycled, converted into biomass, and passed along the food web. In addition, the release of DOM from lysed host cells by viruses contributes to organic particle generation. However, the role of the virioneuston for the microbial loop has never been investigated.
Measurement
Devices used to sample the concentrations of particulates and compounds of the SML include a glass fabric, metal mesh screens, and other hydrophobic surfaces. These are placed on a rotating cylinder which collects surface samples as it rotates on top of the ocean surface.
The glass plate sampler is commonly used. It was first described in 1972 by Harvey and Burzell as a simple but effective method of collecting small sea surface microlayer samples. A clean glass plate is immersed vertically into the water and then withdrawn in a controlled manner. Harvey and Burzell used a plate which was 20 cm square and 4 mm thick. They withdrew it from the sea at the rate of 20 cm per second. Typically the uppermost 20–150 μm of the surface microlayer adheres to the plate as it is withdrawn. The sample is then wiped from both sides of the plate into a sampling vial.
For a plate of the size used by Harvey and Burzel, the resulting sample volumes are between about 3 and 12 cubic centimetres. The sampled SML thickness h in micrometres is given by:
where V is the sample volume in cm3, A is the total immersed plate area of both sides in cm2, and N is the number of times the sample was dipped.
Remote sensing
Ocean surface habitats sit at the interface between the ocean and the atmosphere. The biofilm-like habitat at the surface of the ocean harbours surface-dwelling microorganisms, commonly referred to as neuston. This vast air–water interface sits at the intersection of major air–water exchange processes spanning more than 70% of the global surface area . Bacteria in the surface microlayer of the ocean, called bacterioneuston, are of interest due to practical applications such as air-sea gas exchange of greenhouse gases, production of climate-active marine aerosols, and remote sensing of the ocean. Of specific interest is the production and degradation of surfactants (surface active materials) via microbial biochemical processes. Major sources of surfactants in the open ocean include phytoplankton, terrestrial runoff, and deposition from the atmosphere.
Unlike coloured algal blooms, surfactant-associated bacteria may not be visible in ocean colour imagery. Having the ability to detect these "invisible" surfactant-associated bacteria using synthetic aperture radar has immense benefits in all-weather conditions, regardless of cloud, fog, or daylight. This is particularly important in very high winds, because these are the conditions when the most intense air-sea gas exchanges and marine aerosol production take place. Therefore, in addition to colour satellite imagery, SAR satellite imagery may provide additional insights into a global picture of biophysical processes at the boundary between the ocean and atmosphere, air-sea greenhouse gas exchanges and production of climate-active marine aerosols.
Aeroplankton
A stream of airborne microorganisms, including marine viruses, bacteria and protists, circles the planet above weather systems but below commercial air lanes. Some peripatetic microorganisms are swept up from terrestrial dust storms, but most originate from marine microorganisms in sea spray. In 2018, scientists reported that hundreds of millions of these viruses and tens of millions of bacteria are deposited daily on every square meter around the planet.
Compared to the sub-surface waters, the sea surface microlayer contains elevated concentration of bacteria and viruses, as well as toxic metals and organic pollutants. These materials can be transferred from the sea-surface to the atmosphere in the form of wind-generated aqueous aerosols due to their high vapor tension and a process known as volatilisation. When airborne, these microbes can be transported long distances to coastal regions. If they hit land they can have detrimental effects on animals, vegetation and human health. Marine aerosols that contain viruses can travel hundreds of kilometers from their source and remain in liquid form as long as the humidity is high enough (over 70%). These aerosols are able to remain suspended in the atmosphere for about 31 days. Evidence suggests that bacteria can remain viable after being transported inland through aerosols. Some reached as far as 200 meters at 30 meters above sea level. It was also noted that the process which transfers this material to the atmosphere causes further enrichment in both bacteria and viruses in comparison to either the SML or sub-surface waters (up to three orders of magnitude in some locations).
Mathematical modeling
The stagnant film model is a mathematical model used to simulate the sea surface microlayer. It is a kinematic model which can be used to describe how gas exchange from the ocean's surface and the atmosphere reaches equilibrium. The model assumes both the ocean and atmosphere are composed mostly of well-mixed, constantly moving fluid layers with the sea surface microlayer present as a permanent thin-film layer in the middle. Gas exchange occurs by molecular diffusion between the two fluid layers through the sea surface microlayer.
| Physical sciences | Oceanography | Earth science |
12271360 | https://en.wikipedia.org/wiki/Locustellidae | Locustellidae | Locustellidae is a newly recognized family of small insectivorous songbirds ("warblers"), formerly placed in the Old World warbler "wastebin" family. It contains the grass warblers, grassbirds, and the Bradypterus "bush warblers". These birds occur mainly in Eurasia, Africa, and the Australian region. The family name is sometimes given as Megaluridae, but Locustellidae has priority.
The species are smallish birds with tails that are usually long and pointed; the scientific name of the genus Megalurus in fact means "the large-tailed one" in plain English. They are less wren-like than the typical shrub-warblers (Cettia), but they are similarly drab brownish or buffy all over. They tend to be larger and slimmer than Cettia though, and many have bold dark streaks on wings and/or underside. Most live in scrubland and frequently hunt food by clambering through thick tangled growth or pursuing it on the ground; they are perhaps the most terrestrial of the "warblers". Very unusual for Passeriformes, the beginning of an evolution towards flightlessness is seen in some taxa.
Among the "warbler and babbler" superfamily Sylvioidea, the Locustellidae are closest to the Malagasy warblers, another newly recognized (and hitherto unnamed) family. The black-capped donacobius (Donacobius atricapillus) is a South American relative derived from the same ancestral stock and not a wren as was long believed.
A comprehensive molecular phylogenetic study of the grassbird family Locustellidae published in 2018 found that many of the genera, as then defined, were non-monophyletic. The resulting revision of the genus level taxonomy involved many changes including the resurrection of the genera Poodytes and Cincloramphus as well as the erection of a new genus Helopsaltes. The former genera Megalurulus and Buettikoferella become junior synonyms of Cincloramphus.
Genera
The family contains 67 species divided into 11 genera.
Robsonius – ground warblers (3 species)
Helopsaltes – grasshopper warblers (6 species)
Locustella – grass warblers (23 species)
Bradypterus – megalurid bush warblers (12 species)
Elaphrornis – Sri Lanka bush warbler
Schoenicola – wide-tailed grassbirds (2 species)
Poodytes – grassbirds (5 species including one extinct)
Malia – malia
Cincloramphus – (12 species)
Megalurus – striated grassbird
Catriscus – fan-tailed grassbird
The relationships between the genera is shown in the following cladogram. It is based on a 2018 study by Per Alström and coworkers.
| Biology and health sciences | Passerida | Animals |
13896574 | https://en.wikipedia.org/wiki/Surgical%20suture | Surgical suture | A surgical suture, also known as a stitch or stitches, is a medical device used to hold body tissues together and approximate wound edges after an injury or surgery. Application generally involves using a needle with an attached length of thread. There are numerous types of suture which differ by needle shape and size as well as thread material and characteristics. Selection of surgical suture should be determined by the characteristics and location of the wound or the specific body tissues being approximated.
In selecting the needle, thread, and suturing technique to use for a specific patient, a medical care provider must consider the tensile strength of the specific suture thread needed to efficiently hold the tissues together depending on the mechanical and shear forces acting on the wound as well as the thickness of the tissue being approximated. One must also consider the elasticity of the thread and ability to adapt to different tissues, as well as the memory of the thread material which lends to ease of use for the operator. Different suture characteristics lend way to differing degrees of tissue reaction and the operator must select a suture that minimizes the tissue reaction while still keeping with appropriate tensile strength.
Needles
Historically, surgeons used reusable needles with holes (called "eyes"), which must be threaded before use just as is done with a needle and thread prior to sewing fabric. The advantage of this is that any combination of thread and needle may be chosen to suit the job at hand. Swaged (or "atraumatic") needles with sutures consist of a pre-packed eyeless needle already attached (by swaging) to a specific length of suture thread. This saves time, and eliminates the most difficult threading of very fine needles and sutures.
Two additional benefits are reduced drag and less potential damage to friable tissue during suturing. In a swaged suture the thread is of narrower diameter than the needle, whereas it protrudes on both sides in an eyed needle. Being narrower, the thread in a swaged suture has less drag when passing through tissue than the needle, and, not protruding, is less likely to traumatize friable tissue, earning the combination the designation "atraumatic".
There are several shapes of surgical needles. These include:
Straight
1/4 circle
3/8 circle
1/2 circle. Subtypes of this needle shape include, from larger to smaller size, CT, CT-1, CT-2 and CT-3.
5/8 circle
compound curve
half curved (also known as ski)
half curved at both ends of a straight segment (also known as canoe)
The ski and canoe needle design allows curved needles to be straight enough to be used in laparoscopic surgery, where instruments are inserted into the abdominal cavity through narrow cannulas.
Needles may also be classified by their point geometry; examples include:
taper (needle body is round and tapers smoothly to a point)
cutting (needle body is triangular and has a sharpened cutting edge on the inside curve)
reverse cutting (cutting edge on the outside)
trocar point or tapercut (needle body is round and tapered, but ends in a small triangular cutting point)
blunt points for sewing friable tissues
side cutting or spatula points (flat on top and bottom with a cutting edge along the front to one side) for eye surgery
Finally, atraumatic needles may be permanently swaged to the suture or may be designed to come off the suture with a sharp straight tug. These "pop-offs" are commonly used for interrupted sutures, where each suture is only passed once and then tied.
Sutures can withstand different amounts of force based on their size; this is quantified by the U.S.P. Needles Pull Specifications.
Thread
Materials
Suture material is often broken down into absorbable thread versus non-absorbable thread, which is further delineated into synthetic fibers versus natural fibers. Another important distinction among suture material is whether it is monofilament or polyfilament (braided)
Monofilament versus polyfilament
Monofilament fibers have less tensile strength but create less tissue trauma and are more appropriate with delicate tissues where tissue trauma can be more significant such as small blood vessels. Polyfilament (braided) sutures are composed of multiple fibers and are generally greater in diameter with greater tensile strength, however, they tend to have greater tissue reaction and theoretically have more propensity to harbor bacteria.
Other properties to consider
Tensile strength: the ability of the suture to hold tissues in place without breaking.
Elasticity: the ability of the suture material to adapt to changing tissues such as in cases of edema.
Tissue reactivity: inflammatory response of the surrounding tissue that can cause materials to break down quicker and lose tensile strength. Non absorbable synthetic suture have the lowest of tissue reactivity, while the absorbable natural fibers have the highest rates of tissue reactivity.
Knot security: the ability of the suture to maintain a knot that holds the thread in place.
Absorbable
Absorbable sutures are either degraded via proteolysis or hydrolysis and should not be utilized on body tissue that would require greater than two months of tensile strength. It is generally used internally during surgery or to avoid further procedures for individuals with low likelihood of returning for suture removal. To-date, the available data indicates that the objective short-term wound outcomes are equivalent for absorbable and non-absorbable sutures, and there is equipoise amongst surgeons.
Natural absorbable
Natural absorbable material includes plain catgut, chromic catgut and fast catgut which are all produced from the collagen extracted from bovine intestines. They are all polyfilaments which have different degradations times ranging from 3–28 days. This material is often used for body tissue with low mechanical or shearing force and rapid healing time.
Plain gut (polyfilament)
Description: Maintains original strength for 7–10 days and full degradation occurs in 10 weeks.
Advantages/disadvantages: Excellent elasticity allowing for adaptation to tissue swelling. Passes through the skin with very little tissue trauma occurrence. Poor handling and high tissue reactivity causing quick loss of tensile strength.
Common use: best used in rapidly healing tissues with good blood supply i.e. mucosal tissues.
Chromic gut (polyfilament)
Description: Maintains original strength for 21–28 days and full degradation occurs in 16–18 weeks.
Advantages/disadvantages: Excellent elasticity allowing for adaptation to tissue swelling. Passes through the skin with very little tissue trauma occurrence. Improved handling and decreased tissue reactivity due to chromic salt coating.
Common use: skin closure (face), mucosa, genitalia.
Fast gut (polyfilament)
Description: Treated with heat to further break down protein and allow for more rapid absorption in bodily tissues. Tensile strength less than a week (3–5 days).
Advantages/disadvantages: Excellent elasticity allowing for adaptation to tissue swelling. Passes through the skin with very little tissue trauma occurrence.
Common use: Advised for skin closure only generally on the mucosa or face.
Synthetic absorbable
Synthetic absorbable material includes polyglactic acid, polyglycolic acid, poliglecaprone, polydioxanone, and polytrimethylene carbonate. Among these are monofilaments, polyfilaments and braided sutures. In general synthetic materials will keep tensile strength for longer due to less local tissue inflammation.
Poliglecaprone (monofilament, Monocryl, Monocryl Plus, Suruglyde)
Description: copolymer of synthetic materials. Loses tensile strength quickly; sixty percent lost in the first week. All strength lost within 3 weeks.
Advantages/disadvantages: high tensile strength, excellent elasticity, excellent cosmetic outcomes, decreased hypertrophic scarring, minimal tissue reaction, good knot security originally; however, the material makes the security unreliable over time, thus it is important to keep ears of material long.
Common use: Advised for subcutaneous and superficial tissue closure.
Polyglycolic acid (polyfilament, Dexon)
Description: synthetic polymer that loses all tensile strength in by 25 days. Either dyed green for visibility or undyed.
Advantages/disadvantages: minimal tissue reaction, good tensile strength, good handling, but poor knot security.
Common use: subcutaneous tissue.
Polyglactin 910 (polyfilament, Vicryl)
Description: loss of all tensile strength in 28 days.
Advantages/disadvantages: minimal tissue reaction, good tensile strength, good knot security,
Common use: subcutaneous tissue, skin closure (avoid dyed Vicryl on face).
Polyglactin 910 Irradiated (polyfilament, Vicryl Rapid)
Description: sourced as vicryl is with irradiation to break down material for quicker absorption. Loss of all tensile strength in 5–7 days.
Advantages/disadvantages: minimal tissue reaction, good tensile strength, fair good handling and good knot security.
Common use: scalp and facial laceration closure.
Polglyconate (monofilament, Maxon)
Description: co polymer product of synthetic materials. Loses 75% of the tensile strength after 40 days.
Advantages/disadvantages: minimal tissue reaction, excellent tensile strength, good handling.
Common use: subcutaneous use often an alternative to PDS due to better handling and slightly superior tensile strength.
Polydioxanone closures (PDS, monofilament)
Description: loss of tensile strength in 36–53 days.
Advantages/disadvantages: minimal tissue reaction, good tensile strength, but poor handling.
Common use: subcutaneous with need of high tensile strength (abdominal incision closure).
Non-absorbable
These sutures hold greater tensile strength for longer periods of time and are not subject to degradation. They are appropriate for tissues with a high degree of mechanical or shear force (tendons, certain skin location). They also supply the operator with greater ease of use due to less thread memory.
Natural
Silk (polyfilament, Permahand, Ethicon; Sofsilk, Covidien)
Description: surgical silk is a protein derived from silkworms that is coated to minimize friction and water absorption.
Advantages/disadvantages: This material has good tensile strength, is easy to handle and has excellent knot security. However, it is rarely used internally due to its significant tissue reaction which causes loss of tensile strength over months.
Common use: Due to advancements in sutures, there is no longer indication for use of surgical silk. However, it is still commonly used in dentistry for mucosal surfaces or to secure surgical tubes on the bodies surface.
Synthetic
Synthetic materials include nylon, polypropylene and surgical steel all of which are monofilaments with great tensile strength.
Nylon (monofilaments, Dermalon, Ethilon)
Description: polyamide
Advantages/disadvantages: Excellent tensile strength. However, poor handling and poor knot security due to high material memory.
Common use: Excellent for superficial skin closure due to minimal tissue reactivity. It is the most commonly used skin suture due to its excellent adaptability to potentially expanding tissues (edema).
Nylon (polyfilaments, Nurolon, Surgilon, Supramid)
Description: polyamide
Advantages/disadvantages: Excellent tensile strength, increased usability, and increased knot security as compared to its monofilamentous counterpart. However, its polyfilamentous nature is said to increase risk of infection.
Common use: soft tissue, vessel ligations and superficial skin (specifically facial lacerations).
Braided polyester (polyfilament, Ethibond, Dagrofil, Synthofil, PremiCron, Synthofil)
Description: made from polyethylene terephthalate, there are various brands and configurations of this type of suture. Many are braided, coated in silicone and dyed for visibility.
Advantages/disadvantages: Good handling, good knot security and high tensile strength due to low tissue reactivity. However, this suture can create more tissue trauma when passing through the skin and is more expensive than its counterparts
Common use: Rare, pediatric valvular surgery, alternative to surgical steel for orthopedic surgery due to superior handling.
Polybutester (monofilament, Novafil)
Description: A copolymer of polyester.
Advantages/disadvantages: low tissue reactivity, good handling, high tensile strength that is greater than most other monofilaments, good elasticity during increasing edema.
Common use: rare, tendon repairs, plastics (pull out subcuticular stitch)
Surgical steel
Description: synthetic mixture of multiple alloys.
Advantages/disadvantages: Tensile strength is exceptional with very little tissue reactivity, thus maintaining minimal degradation over time. This suture material has very poor handling.
Common use: orthopedics, sternum closure.
Sizes
Suture sizes are defined by the United States Pharmacopeia (U.S.P.). Sutures were originally manufactured ranging in size from #1 to #6, with #1 being the smallest. A #4 suture would be roughly the diameter of a tennis racquet string. The manufacturing techniques, derived at the beginning from the production of musical strings, did not allow thinner diameters. As the procedures improved, #0 was added to the suture diameters, and later, thinner and thinner threads were manufactured, which were identified as #00 (#2-0 or #2/0) to #000000 (#6-0 or #6/0).
Modern sutures range from #5 (heavy braided suture for orthopedics) to #11-0 (fine monofilament suture for ophthalmics). Atraumatic needles are manufactured in all shapes for most sizes. The actual diameter of thread for a given U.S.P. size differs depending on the suture material class.
Techniques
Many different techniques exist. The most common is the simple interrupted stitch; it is indeed the simplest to perform and is called "interrupted" because the suture thread is cut between each individual stitch. The vertical and horizontal mattress stitch are also interrupted but are more complex and specialized for everting the skin and distributing tension. The running or continuous stitch is quicker but risks failing if the suture is cut in just one place; the continuous locking stitch is in some ways a more secure version. The chest drain stitch and corner stitch are variations of the horizontal mattress.
Other stitches or suturing techniques include:
Purse-string suture, a continuous, circular inverting suture which is made to secure apposition of the edges of a surgical or traumatic wound.
Figure-of-eight stitch
Subcuticular stitch. A continuous suture where the needle enters and exits the epidermis along the plane of the skin. This stitch is for approximating superficial skin edges and provides the best cosmetic result. Superficial gapping wounds may be reduced effectively by using continuous subcuticular sutures. It is unclear whether subcuticular sutures can reduce the rate of surgical site infections.when compared with other suturing methods.
Placement
Sutures are placed by mounting a needle with attached suture into a needle holder. The needle point is pressed into the flesh, advanced along the trajectory of the needle's curve until it emerges, and pulled through. The trailing thread is then tied into a knot, usually a square knot or surgeon's knot. Ideally, sutures bring together the wound edges, without causing indenting or blanching of the skin, since the blood supply may be impeded and thus increase infection and scarring. Ideally, sutured skin rolls slightly outward from the wound (eversion), and the depth and width of the sutured flesh is roughly equal. Placement varies based on the location,
Stitching interval and spacing
Skin and other soft tissue can lengthen significantly under strain. To accommodate this lengthening, continuous stitches must have an adequate amount of slack. Jenkin's rule was the first research result in this area, showing that the then-typical use of a suture-length to wound-length ratio of 2:1 increased the risk of a burst wound, and suggesting a SL:WL ratio of 4:1 or more in abdominal wounds. A later study suggested 6:1 as the optimal ratio in abdominal closure.
Layers
In contrast to single layer suturing, two layer suturing generally involves suturing at a deeper level of a tissue followed by another layer of suturing at a more superficial level. For example, Cesarean section can be performed with single or double layer suturing of the uterine incision.
Removal
Whereas some sutures are intended to be permanent, and others in specialized cases may be kept in place for an extended period of many weeks, as a rule sutures are a short-term device to allow healing of a trauma or wound.
Removal of sutures is traditionally achieved by using forceps to hold the suture thread steady and pointed scalpel blades or scissors to cut. For practical reasons the two instruments (forceps and scissors) are available in a sterile kit. In certain countries (e.g. US), these kits are available in sterile disposable trays because of the high cost of cleaning and re-sterilization.
Expansions
A pledgeted suture is one that is supported by a pledget, that is, a small flat non-absorbent pad normally composed of polytetrafluoroethylene, used as buttresses under sutures when there is a possibility of sutures tearing through tissue.
Tissue adhesives
Topical cyanoacrylate adhesives (closely related to super glue), have been used in combination with, or as an alternative to, sutures in wound closure. The adhesive remains liquid until exposed to water or water-containing substances/tissue, after which it cures (polymerizes) and forms a bond to the underlying surface. The tissue adhesive has been shown to act as a barrier to microbial penetration as long as the adhesive film remains intact. Limitations of tissue adhesives include contraindications to use near the eyes and a mild learning curve on correct usage. They are also unsuitable for oozing or potentially contaminated wounds.
In surgical incisions it does not work as well as sutures as the wounds often break open.
Cyanoacrylate is the generic name for cyanoacrylate based fast-acting glues such as methyl-2-cyanoacrylate, ethyl-2-cyanoacrylate (commonly sold under trade names like Superglue and Krazy Glue) and n-butyl-cyanoacrylate. Skin glues like Indermil and Histoacryl were the first medical grade tissue adhesives to be used, and these are composed of n-butyl cyanoacrylate. These worked well but had the disadvantage of having to be stored in the refrigerator, were exothermic so they stung the patient, and the bond was brittle. Nowadays, the longer chain polymer, 2-octyl cyanoacrylate, is the preferred medical grade glue. It is available under various trade names, such as LiquiBand, SurgiSeal, FloraSeal, and Dermabond. These have the advantages of being more flexible, making a stronger bond, and being easier to use. The longer side chain types, for example octyl and butyl forms, also reduce tissue reaction.
History
Through many millennia, various suture materials were used or proposed. Needles were made of bone or metals such as silver, copper, and aluminium bronze wire. Sutures were made of plant materials (flax, hemp and cotton) or animal material (hair, tendons, arteries, muscle strips and nerves, silk, and catgut).
The earliest reports of surgical suture date to 3000 BC in ancient Egypt, and the oldest known suture is in a mummy from 1100 BC. A detailed description of a wound suture and the suture materials used in it is by the Indian sage and physician Sushruta, written in 500 BC. The Greek father of medicine, Hippocrates, described suture techniques, as did the later Roman Aulus Cornelius Celsus. The 2nd-century Roman physician Galen described sutures made of surgical gut or catgut. In the 10th century, the catgut suture along with the surgery needle were used in operations by Abulcasis. The gut suture was similar to that of strings for violins, guitars, and tennis racquets and it involved harvesting sheep or cow intestines. Catgut sometimes led to infection due to a lack of disinfection and sterilization of the material.
Joseph Lister endorsed the routine sterilization of all suture threads. He first attempted sterilization with the 1860s "carbolic catgut", and chromic catgut followed two decades later. Sterile catgut was finally achieved in 1906 with iodine treatment.
The next great leap came in the twentieth century. The chemical industry drove production of the first synthetic thread in the early 1930s, which exploded into production of numerous absorbable and non-absorbable synthetics. The first synthetic absorbable was based on polyvinyl alcohol in 1931. Polyesters were developed in the 1950s, and later the process of radiation sterilization was established for catgut and polyester. Polyglycolic acid was discovered in the 1960s and implemented in the 1970s. Today, most sutures are made of synthetic polymer fibers. Silk and, rarely, gut sutures are the only materials still in use from ancient times. In fact, gut sutures have been banned in Europe and Japan owing to concerns regarding bovine spongiform encephalopathy. Silk suture is still used today, mainly to secure surgical drains.
| Technology | Surgical instruments | null |
14990054 | https://en.wikipedia.org/wiki/Sleep%20in%20animals | Sleep in animals | Sleep appears to be a biological requirement for all animals except for basal species with no brain or only a rudimentary brain. It has been observed in mammals, birds, reptiles, amphibians, fish, and, in some form, in insects. The internal circadian clock promotes sleep at night for diurnal organisms (such as humans) and in the day for nocturnal organisms (such as rats). Sleep patterns vary widely among species, with some foregoing sleep for extended periods and some engaging in unihemispheric sleep, in which one brain hemisphere sleeps while the other remains awake.
Definition
Sleep can follow a physiological or behavioral definition. In the physiological sense, sleep is a state characterized by reversible unconsciousness, special brainwave patterns, sporadic eye movement, loss of muscle tone (possibly with some exceptions; see below regarding the sleep of birds and of aquatic mammals), and a compensatory increase following deprivation of the state, this last known as sleep homeostasis (i.e., the longer a waking state lasts, the greater the intensity and duration of the sleep state thereafter). In the behavioral sense, sleep is characterized by minimal movement, non-responsiveness to external stimuli (i.e. increased sensory threshold), the adoption of a typical posture, and the occupation of a sheltered site, all of which is usually repeated on a 24-hour basis. The physiological definition applies well to birds and mammals, but in other animals whose brains are not as complex, the behavioral definition is more often used. In very simple animals, behavioral definitions of sleep are the only ones possible, and even then the behavioral repertoire of the animal may not be extensive enough to allow distinction between sleep and wakefulness. Sleep is quickly reversible, as opposed to hibernation or coma, and sleep deprivation is followed by longer or deeper rebound sleep.
Necessity
If sleep were not essential, one would expect to find
Animal species that do not sleep at all
Animals that do not need recovery sleep after staying awake longer than usual
Animals that suffer no serious consequences as a result of lack of sleep
These symptoms are not seen in complex animals, and sleep is thus considered necessary to them. Sleep helps the body and mind to feel rested. Findings show that if rats do not get sleep, they die in a few weeks. Despite having enough food, their appetite tends to decrease resulting in weight loss and eventually death.
Outside of a few basal animals that have no brain or a very simple one, no animals have been found to date that satisfy any of these criteria. While some varieties of shark, such as great whites and hammerheads, must remain in motion at all times to move oxygenated water over their gills, it is possible they still sleep one cerebral hemisphere at a time as marine mammals do. However, it remains to be shown definitively whether any fish is capable of unihemispheric sleep.
Invertebrates
Sleep as a phenomenon appears to have very old evolutionary roots. Unicellular organisms do not necessarily "sleep", although many of them have pronounced circadian rhythms. The fresh-water polyp Hydra vulgaris and the jellyfish Cassiopea are among the most primitive organisms in which sleep-like states have been observed. Observing sleep states in jellyfish provides evidence that sleep states do not require that an animal have a brain or central nervous system. The nematode C. elegans is another primitive organism that appears to require sleep. Here, a lethargus phase occurs in short periods preceding each moult, a fact which may indicate that sleep primitively is connected to developmental processes. Raizen et al.'''s results furthermore suggest that sleep is necessary for changes in the neural system.
Bees have some of the most complex sleep states amongst insects.
Decade after decade results mounted that insects do sleep, and that this resembles mammalian and avian sleep. Nonetheless, sleep scientists continued to not accept these results and there was wide agreement that insects did not experience sleep. It took the gene expression studies of Hendricks et al. 2000 and Shaw et al. 2000 showing orthology between mammals and the fruit fly Drosophila melanogaster for this to finally be accepted. The electrophysiological study of sleep in small invertebrates is complicated. Insects go through circadian rhythms of activity and passivity but some do not seem to have a homeostatic sleep need. Insects do not seem to exhibit REM sleep. However, fruit flies appear to sleep, and systematic disturbance of that state leads to cognitive disabilities. There are several methods of measuring cognitive functions in fruit flies. A common method is to let the flies choose whether they want to fly through a tunnel that leads to a light source, or through a dark tunnel. Normally, flies are attracted to light. But if sugar is placed in the end of the dark tunnel, and something the flies dislike is placed in the end of the light tunnel, the flies will eventually learn to fly towards darkness rather than light. Flies deprived of sleep require a longer time to learn this and also forget it more quickly. If an arthropod is experimentally kept awake longer than it is used to, then its coming rest period will be prolonged. In cockroaches, that rest period is characterized by the antennae being folded down and by a decreased sensitivity to external stimuli. Sleep has been described in crayfish, too, characterized by passivity and increased thresholds for sensory stimuli as well as changes in the EEG pattern, markedly differing from the patterns found in crayfish when they are awake. In honeybees, it has been shown that they use sleep to store long-term memories. Sleep-like state has been described in jumping spiders, too, as well as regularly occurring bouts of retinal movements that suggest an
REM sleep–like state. Also sleeping cuttlefish and octopuses show signs of having REM-sleep behaviors.
Fish
Sleep in fish is the subject of ongoing scientific research.Reebs, S. (1992) Sleep, inactivity and circadian rhythms in fish. pp. 127–135 in: Ali, M.A. (ed.), Rhythms in Fish, New York: Plenum Press. Typically fish exhibit periods of inactivity but show no significant reactions to deprivation of this condition. Some species that always live in shoals or that swim continuously (because of a need for ram ventilation of the gills, for example) are suspected never to sleep. There is also doubt about certain blind species that live in caves. Other fish seem to sleep, however. For example, zebrafish,; third party discussion of Yokogawa: tilapia, tench, brown bullhead, and swell shark become motionless and unresponsive at night (or by day, in the case of the swell shark); Spanish hogfish and blue-headed wrasse can even be lifted by hand all the way to the surface without evoking a response. Studies show that some fish (for example rays and sharks) have unihemispheric sleep, which means they put half their brain to sleep while the other half still remains active and they swim while they are sleeping. A 1961 observational study of approximately 200 species in European public aquaria reported many cases of apparent sleep. On the other hand, sleep patterns are easily disrupted and may even disappear during periods of migration, spawning, and parental care.
Land vertebrates
Mammals, birds and reptiles evolved from amniotic ancestors, the first vertebrates with life cycles independent of water. The fact that birds and mammals are the only known animals to exhibit REM and NREM sleep indicates a common trait before divergence. However, recent evidence of REM-like sleep in fish suggests this divergence may have occurred much earlier than previously thought. Up to this point, reptiles were considered the most logical group to investigate the origins of sleep. Daytime activity in reptiles alternates between basking and short bouts of active behavior, which has significant neurological and physiological similarities to sleep states in mammals. It is proposed that REM sleep evolved from short bouts of motor activity in reptiles, while slow-wave sleep (SWS) evolved from their basking state, which shows similar slow -wave EEG patterns.
Reptiles and amphibians
Reptiles have quiescent periods similar to mammalian sleep, and a decrease in electrical activity in the brain has been registered when the animals have been asleep. However, the EEG pattern in reptilian sleep differs from what is seen in mammals and other animals. In reptiles, sleep time increases following sleep deprivation, and stronger stimuli are needed to awaken the animals when they have been deprived of sleep as compared to when they have slept normally. This suggests that the sleep which follows deprivation is compensatorily deeper.
In 2016, a study reported the existence of REM- and NREM-like sleep stages in the Australian dragon Pogona vitticeps. Amphibians have periods of inactivity but show high vigilance (receptivity to potentially threatening stimuli) in this state.
Like some birds and aquatic mammals, crocodilians are also capable of unihemispheric sleep.
Birds
There are significant similarities between sleep in birds and sleep in mammals, which is one of the reasons for the idea that sleep in higher animals with its division into REM and NREM sleep has evolved together with warm-bloodedness. Birds compensate for sleep loss in a manner similar to mammals, by deeper or more intense slow-wave sleep (SWS).
Birds have both REM and NREM sleep, and the EEG patterns of both have similarities to those of mammals. Different birds sleep different amounts, but the associations seen in mammals between sleep and variables such as body mass, brain mass, relative brain mass, basal metabolism and other factors (see below) are not found in birds. The only clear explanatory factor for the variations in sleep amounts for birds of different species is that birds who sleep in environments where they are exposed to predators have less deep sleep than birds sleeping in more protected environments.
Birds do not necessarily exhibit sleep debt, but a peculiarity that birds share with aquatic mammals, and possibly also with certain species of lizards (opinions differ about that last point), is the phenomenon of unihemispheric slow-wave sleep; that is, the ability to sleep with one cerebral hemisphere at a time, while keeping the other hemisphere awake. When just one hemisphere is sleeping, only the contralateral eye will be shut; that is, when the right hemisphere is asleep, the left eye will be shut, and vice versa. The distribution of sleep between the two hemispheres and the amount of unihemispheric sleep are determined both by which part of the brain has been the most active during the previous period of wake—that part will sleep the deepest—and by the level of risk of attacks from predators. Ducks near the perimeter of the flock are likely to be the ones that first will detect predator attacks. These ducks have significantly more unihemispheric sleep than those who sleep in the middle of the flock, and they react to threatening stimuli seen by the open eye.
Opinions partly differ about sleep in migratory birds. The controversy is mainly about whether they can sleep while flying or not. Theoretically, certain types of sleep could be possible while flying, but technical difficulties preclude the recording of brain activity in birds while they are flying.
Mammals
Mammals have wide diversity in sleep phenomena. Generally, they go through periods of alternating non-REM and REM sleep, but these manifest differently. Horses and other herbivorous ungulates can sleep while standing, but must necessarily lie down for REM sleep (which causes muscular atony) for short periods. Giraffes, for example, only need to lie down for REM sleep for a few minutes at a time. Bats sleep while hanging upside down. Male armadillos get erections during non-REM sleep, and the inverse is true in rats. Early mammals engaged in polyphasic sleep, dividing sleep into multiple bouts per day. Higher daily sleep quotas and shorter sleep cycles in polyphasic species as compared to monophasic species, suggest that polyphasic sleep may be a less efficient means of attaining sleep's benefits. Small species with higher basal metabolic rate (BMR) may therefore have less efficient sleep patterns. It follows that the evolution of monophasic sleep may hitherto be an unknown advantage of evolving larger mammalian body sizes and therefore lower BMR.
Sleep is sometimes thought to help conserve energy, though this theory is not fully adequate as it only decreases metabolism by about 5–10%."Function of Sleep.". Scribd.com. Retrieved on 1 December 2011. Additionally it is observed that mammals require sleep even during the hypometabolic state of hibernation, in which circumstance it is actually a net loss of energy as the animal returns from hypothermia to euthermia in order to sleep.
Nocturnal animals have higher body temperatures, greater activity, rising serotonin, and diminishing cortisol during the night—the inverse of diurnal animals. Nocturnal and diurnal animals both have increased electrical activity in the suprachiasmatic nucleus, and corresponding secretion of melatonin from the pineal gland, at night. Nocturnal mammals, which tend to stay awake at night, have higher melatonin at night just like diurnal mammals do. And, although removing the pineal gland in many animals abolishes melatonin rhythms, it does not stop circadian rhythms altogether—though it may alter them and weaken their responsiveness to light cues. Cortisol levels in diurnal animals typically rise throughout the night, peak in the awakening hours, and diminish during the day.Thomas A. Wehr (1999). "The Impact of Changes in Nightlength (Scotoperiod) on Human Sleep", in Turek & Zee (eds.), Regulation of Sleep and Circadian Rhythms, pp. 263–285. In diurnal animals, sleepiness increases during the night.
Duration
Different mammals sleep different amounts. Some, such as bats, sleep 18–20 hours per day, while others, including giraffes, sleep only 3–4 hours per day. There can be big differences even between closely related species. There can also be big differences between laboratory and field studies: for example, researchers in 1983 reported that captive sloths slept nearly 16 hours a day, but in 2008, when miniature neurophysiological recorders were developed that could be affixed to wild animals, sloths in nature were found to sleep only 9.6 hours a day.
As with birds, the main rule for mammals (with certain exceptions, see below) is that they have two essentially different stages of sleep: REM and NREM sleep (see above). Mammals' feeding habits are associated with their sleep length. The daily need for sleep is highest in carnivores, lower in omnivores and lowest in herbivores. Humans sleep less than many other omnivores but otherwise not unusually much or unusually little in comparison with other mammals.
Many herbivores, like Ruminantia (such as cattle), spend much of their wake time in a state of drowsiness, which perhaps could partly explain their relatively low need for sleep. In herbivores, an inverse correlation is apparent between body mass and sleep length; big mammals sleep less than smaller ones. This correlation is thought to explain about 25% of the difference in sleep amount between different mammals. Also, the length of a particular sleep cycle is associated with the size of the animal; on average, bigger animals will have sleep cycles of longer durations than smaller animals. Sleep amount is also coupled to factors like basal metabolism, brain mass, and relative brain mass. The duration of sleep among species is also directly related to BMR. Rats, which have a high BMR, sleep for up to 14 hours a day, whereas elephants and giraffes, which have lower BMRs, sleep only 2–4 hours per day.
It has been suggested that mammalian species which invest in longer sleep times are investing in the immune system, as species with the longer sleep times have higher white blood cell counts. Mammals born with well-developed regulatory systems, such as the horse and giraffe, tend to have less REM sleep than the species which are less developed at birth, such as cats and rats. This appears to echo the greater need for REM sleep among newborns than among adults in most mammal species. Many mammals sleep for a large proportion of each 24-hour period when they are very young. The giraffe only sleeps 2 hours a day in about 5–15 minute sessions. Koalas are the longest sleeping-mammals, about 20–22 hours a day. However, killer whales and some other dolphins do not sleep during the first month of life. Instead, young dolphins and whales frequently take rests by pressing their body next to their mother's while she swims. As the mother swims she is keeping her offspring afloat to prevent them from drowning. This allows young dolphins and whales to rest, which will help keep their immune system healthy; in turn, protecting them from illnesses. During this period, mothers often sacrifice sleep for the protection of their young from predators. However, unlike other mammals, adult dolphins and whales are able to go without sleep for a month.
Comparative average sleep periods for various mammals (in captivity) over 24 hours
Horses – 2 hours
Elephants – 3+ hours
Cows – 4.0 hours
Giraffes – 4.5 hours
Humans – 8.0 hours
Rabbits – 8.4 hours
Chimpanzees – 9.7 hours
Red foxes – 9.8 hours
Dogs – 10.1 hours
House mice – 12.5 hours
Cats – 12.5 hours
Lions – 13.5 hours
Platypuses – 14 hours
Chipmunks – 15 hours
Tigers – 15.8 hours
Giant armadillos – 18.1 hours
Leopards – 18 hours
Little brown bats – 19.9 hours
Reasons given for the wide variations include the fact that mammals "that nap in hiding, like bats or rodents tend to have longer, deeper snoozes than those on constant alert." Lions, which have little fear of predators also have relatively long sleep periods, while elephants have to eat most of the time to support their huge bodies. Little brown bats conserve their energy except for the few hours each night when their insect prey are available, and platypuses eat a high energy crustacean diet and, therefore, probably do not need to spend as much time awake as many other mammals.
Rodents
A study conducted by Datta indirectly supports the idea that memory benefits from sleep. A box was constructed wherein a single rat could move freely from one end to the other. The bottom of the box was made of a steel grate. A light would shine in the box accompanied by a sound. After a five-second delay, an electrical shock would be applied. Once the shock commenced, the rat could move to the other end of the box, ending the shock immediately. The rat could also use the five-second delay to move to the other end of the box and avoid the shock entirely. The length of the shock never exceeded five seconds. This was repeated 30 times for half the rats. The other half, the control group, was placed in the same trial, but the rats were shocked regardless of their reaction. After each of the training sessions, the rat would be placed in a recording cage for six hours of polygraphic recordings. This process was repeated for three consecutive days. During the posttrial sleep recording session, rats spent 25.47% more time in REM sleep after learning trials than after control trials.
An observation of the Datta study is that the learning group spent 180% more time in SWS than did the control group during the post-trial sleep-recording session. This study shows that after spatial exploration activity, patterns of hippocampal place cells are reactivated during SWS following the experiment. Rats were run through a linear track using rewards on either end. The rats would then be placed in the track for 30 minutes to allow them to adjust (PRE), then they ran the track with reward-based training for 30 minutes (RUN), and then they were allowed to rest for 30 minutes.
During each of these three periods, EEG data were collected for information on the rats' sleep stages. The mean firing rates of hippocampal place cells during prebehavior SWS (PRE) and three ten-minute intervals in postbehavior SWS (POST) were calculated by averaging across 22 track-running sessions from seven rats. The results showed that ten minutes after the trial RUN session, there was a 12% increase in the mean firing rate of hippocampal place cells from the PRE level. After 20 minutes, the mean firing rate returned rapidly toward the PRE level. The elevated firing of hippocampal place cells during SWS after spatial exploration could explain why there were elevated levels of slow-wave sleep in Datta's study, as it also dealt with a form of spatial exploration.
In rats, sleep deprivation causes weight loss and reduced body temperature. Rats kept awake indefinitely develop skin lesions, hyperphagia, loss of body mass, hypothermia, and, eventually, fatal sepsis. Sleep deprivation also hinders the healing of burns on rats. When compared with a control group, sleep-deprived rats' blood tests indicated a 20% decrease in white blood cell count, a significant change in the immune system.
A 2014 study found that depriving mice of sleep increased cancer growth and dampened the immune system's ability to control cancers. The researchers found higher levels of M2 tumor-associated macrophages and TLR4 molecules in the sleep deprived mice and proposed this as the mechanism for increased susceptibility of the mice to cancer growth. M2 cells suppress the immune system and encourage tumour growth. TRL4 molecules are signalling molecules in the activation of the immune system.
Monotremes
Since monotremes (egg-laying mammals) are considered to represent one of the evolutionarily oldest groups of mammals, they have been subject to special interest in the study of mammalian sleep. As early studies of these animals could not find clear evidence for REM sleep, it was initially assumed that such sleep did not exist in monotremes, but developed after the monotremes branched off from the rest of the mammalian evolutionary line, and became a separate, distinct group. However, EEG recordings of the brain stem in monotremes show a firing pattern that is quite similar to the patterns seen in REM sleep in higher mammals. In fact, the largest amount of REM sleep known in any animal is found in the platypus. REM electrical activation does not extend at all to the forebrain in platypods, suggesting that they do not dream. The average sleep time of the platypus in a 24-hour period is said to be as long as 14 hours, though this may be because of their high-calorie crustacean diet.
Aquatic mammals
The consequences of falling into a deep sleep for marine mammalian species can be suffocation and drowning, or becoming easy prey for predators. Thus, dolphins, whales, and pinnipeds (seals) engage in unihemispheric sleep while swimming, which allows one brain hemisphere to remain fully functional, while the other goes to sleep. The hemisphere that is asleep alternates, so that both hemispheres can be fully rested. Just like terrestrial mammals, pinnipeds that sleep on land fall into a deep sleep and both hemispheres of their brain shut down and are in full sleep mode. Aquatic mammal infants do not have REM sleep in infancy; REM sleep increases as they age.
Among others, seals and whales belong to the aquatic mammals. Earless seals and eared seals have solved the problem of sleeping in water via two different methods. Eared seals, like whales, show unihemispheric sleep. The sleeping half of the brain does not awaken when they surface to breathe. When one half of a seal's brain shows slow-wave sleep, the flippers and whiskers on its opposite side are immobile. While in the water, these seals have almost no REM sleep and may go a week or two without it. As soon as they move onto land they switch to bilateral REM sleep and NREM sleep comparable to land mammals, surprising researchers with their lack of "recovery sleep" after missing so much REM.
Earless seals sleep bihemispherically like most mammals, under water, hanging at the water surface or on land. They hold their breath while sleeping under water, and wake up regularly to surface and breathe. They can also hang with their nostrils above water and in that position have REM sleep, but they do not have REM sleep underwater.
REM sleep has been observed in the pilot whale, a species of dolphin. Whales do not seem to have REM sleep, nor do they seem to have any problems because of this. One reason REM sleep might be difficult in marine settings is the fact that REM sleep causes muscular atony; that is to say, a functional paralysis of skeletal muscles that can be difficult to combine with the need to breathe regularly.
Unihemispherism
Unihemispheric sleep refers to sleeping with only a single cerebral hemisphere. The phenomenon has been observed in birds and aquatic mammals, as well as in several reptilian species (the latter being disputed: many reptiles behave in a way which could be construed as unihemispheric sleeping, but EEG studies have given contradictory results). Reasons for the development of unihemispheric sleep are likely that it enables the sleeping animal to receive stimuli—threats, for instance—from its environment, and that it enables the animal to fly or periodically surface to breathe when immersed in water. Only NREM sleep exists unihemispherically, and there seems to exist a continuum in unihemispheric sleep regarding the differences in the hemispheres: in animals exhibiting unihemispheric sleep, conditions range from one hemisphere being in deep sleep with the other hemisphere being awake to one hemisphere sleeping lightly with the other hemisphere being awake. If one hemisphere is selectively deprived of sleep in an animal exhibiting unihemispheric sleep (one hemisphere is allowed to sleep freely but the other is awoken whenever it falls asleep), the amount of deep sleep will selectively increase in the hemisphere that was deprived of sleep when both hemispheres are allowed to sleep freely.
The neurobiological background for unihemispheric sleep is still unclear. In experiments on cats in which the connection between the left and the right halves of the brain stem has been severed, the brain hemispheres show periods of a desynchronized EEG, during which the two hemispheres can sleep independently of each other. In these cats, the state where one hemisphere slept NREM and the other was awake, as well as one hemisphere sleeping NREM with the other state sleeping REM were observed. The cats were never seen to sleep REM sleep with one hemisphere while the other hemisphere was awake. This is in accordance with the fact that REM sleep, as far as is currently known, does not occur unihemispherically.
The fact that unihemispheric sleep exists has been used as an argument for the necessity of sleep. It appears that no animal has developed an ability to go without sleep altogether.
Hibernation
Animals that hibernate are in a state of torpor, differing from sleep. Hibernation markedly reduces the need for sleep, but does not remove it. Some hibernating animals end their hibernation a couple of times during the winter so that they can sleep. Hibernating animals waking up from hibernation often go into rebound sleep because of lack of sleep during the hibernation period. They are definitely well-rested and are conserving energy during hibernation, but need sleep for something else.
Dreams
Dreaming in dogs has been studied by Stanley Coren, Professor Emeritus of Psychology at the University of British Columbia in Vancouver. Researchers have studied dreaming in dogs by manipulating the pons in the brain stem. He is the author of the book Do Dogs Dream? Nearly Everything Your Dog Wants You to Know.'' (Norton, 2012).
| Biology and health sciences | Ethology | Biology |
3239188 | https://en.wikipedia.org/wiki/ARINC%20429 | ARINC 429 | ARINC 429, the "Mark 33 Digital Information Transfer System (DITS)," is the ARINC technical standard for the predominant avionics data bus used on most higher-end commercial and transport aircraft. It defines the physical and electrical interfaces of a two-wire data bus and a data protocol to support an aircraft's avionics local area network.
Technical description
Medium and signaling
ARINC 429 is a data transfer standard for aircraft avionics. It uses a self-clocking, self-synchronizing data bus protocol (Tx and Rx are on separate ports). The physical connection wires are twisted pairs carrying balanced differential signaling. Data words are 32 bits in length and most messages consist of a single data word. Messages are transmitted at either 12.5 or 100 kbit/s to other system elements that are monitoring the bus messages. The transmitter constantly transmits either 32-bit data words or the NULL state (0 Volts). A single wire pair is limited to one transmitter and no more than 20 receivers. The protocol allows for self-clocking at the receiver end, thus eliminating the need to transmit clocking data. ARINC 429 is an alternative to MIL-STD-1553.
Bit numbering, transmission order, and bit significance
The ARINC 429 unit of transmission is a fixed-length 32-bit frame, which the standard refers to as a 'word'. The bits within an ARINC 429 word are serially identified from Bit Number 1 to Bit Number 32 or simply Bit 1 to Bit 32. The fields and data structures of the ARINC 429 word are defined in terms of this numbering.
While it is common to illustrate serial protocol frames progressing in time from right to left, a reversed ordering is commonly practiced within the ARINC standard. Even though ARINC 429 word transmission begins with Bit 1 and ends with Bit 32, it is common to diagram and describe ARINC 429 words in the order from Bit 32 to Bit 1.
In simplest terms, while the transmission order of bits (from the first transmitted bit to the last transmitted bit) for a 32-bit frame is conventionally diagrammed as
First bit > 1, 2, 3, 4, 5, 6, 7, 8, 9, 10, 11, 12, ... 29, 30, 31, 32 < Last bit,
this sequence is often diagrammed in ARINC 429 publications in the opposite direction as
Last bit > 32, 31, 30, 29, ... 12, 11, 10, 9, 8, 7, 6, 5, 4, 3, 2, 1 < First bit.
Generally, when the ARINC 429 word format is illustrated with Bit 32 to the left, the numeric representations in the data field are read with the most significant bit on the left. However, in this particular bit order presentation, the Label field reads with its most significant bit on the right. Like CAN Protocol Identifier Fields, ARINC 429 label fields are transmitted most significant bit first. However, like UART Protocol, Binary-coded decimal numbers and binary numbers in the ARINC 429 data fields are generally transmitted least significant bit first.
Some equipment suppliers publish the bit transmission order as
First bit > 8, 7, 6, 5, 4, 3, 2, 1, 9, 10, 11, 12, 13 … 32 < Last bit.
The suppliers that use this representation have in effect renumbered the bits in the Label field, converting the standard's MSB 1 bit numbering for that field to LSB 1 bit numbering. This renumbering highlights the relative reversal of "bit endianness" between the Label representation and numeric data representations as defined within the ARINC 429 standard. Of note is how the 87654321 bit numbering is similar to the 76543210 bit numbering common in digital equipment; but reversed from the 12345678 bit numbering defined for the ARINC 429 Label field.
This notional reversal also reflects historical implementation details. ARINC 429 transceivers have been implemented with 32-bit shift registers. Parallel access to that shift register is often octet-oriented. As such, the bit order of the octet access is the bit order of the accessing device, which is usually LSB 0; and serial transmission is arranged such that the least significant bit of each octet is transmitted first. So, in common practice, the accessing device wrote or read a "reversed label" (for example, to transmit a Label 2138 [or 8B16] the bit-reversed value D116 is written to the Label octet). Newer or "enhanced" transceivers may be configured to reverse the Label field bit order "in hardware."
Word format
Each ARINC 429 word is a 32-bit sequence that contains five fields:
Bit 32 is the parity bit, and is used to verify that the word was not damaged or garbled during transmission. Every ARINC 429 channel typically uses "odd" parity - there must be an odd number of "1" bits in the word. This bit is set to 0 or 1 to ensure that the correct number of bits are set to 1 in the word.
Bits 30 to 31 are the Sign/Status Matrix (SSM) - these bits may have various encodings dependent on the particular data representation applied to a given word:
In all cases using the SSM, these bits may be encoded to indicate:
Normal Operation (NO) - Indicates the data in this word is considered to be correct data.
Functional Test (FT) - Indicates that the data is being provided by a test source.
Failure Warning (FW) - Indicates a failure which causes the data to be suspect or missing.
No Computed Data (NCD) - Indicates that the data is missing or inaccurate for some reason other than a failure. For example, autopilot commands will show as NCD when the autopilot is not turned on.
In the case of Binary Coded Decimal (BCD) representation, the SSM may also indicate the Sign (+/-) of the data or some information analogous to sign, like an orientation (North/South; East/West). When so indicating sign, the SSM is also considered to be indicating Normal Operation.
In the case of two's-complement representation of signed binary numbers (BNR), Bit 29 represents the number's sign; that is, sign indication is delegated to Bit 29 in this case.
In the case of discrete data representation (e.g., bit-fields), the SSM has a different, signless encoding.
{| class="wikitable" border="1"
|-
! colspan="2" align="center" | SSM
! colspan="3" align="center" | Data Dependent SSM Encodings:
|-
! align="left" | Bit 31
! Bit 30
! Sign/Status Matrix for BCD Data
! Status Matrix for BNR Data
! Status Matrix for Discrete Data
|-
| align="center" | 0
| align="center" | 0
| Plus, North, East, Right, To, Above
| Failure Warning (FW)
| Verified Data, Normal Operation
|-
| align="center" | 0
| align="center" | 1
| No Computed Data (NCD)
| No Computed Data (NCD)
| No Computed Data (NCD)
|-
| align="center" | 1
| align="center" | 0
| Functional Test (FT)
| Functional Test (FT)
| Functional Test (FT)
|-
| align="center" | 1
| align="center" | 1
| Minus, South, West, Left, From, Below
| Normal Operation (NO)
| Failure Warning (FW)
|-
|}
{| class="wikitable" border="1"
|-
! Bit 29
! Sign Matrix for BNR Data
|-
| align="center" | 0
| Plus, North, East, Right, To, Above
|-
| align="center" | 1
| Minus, South, West, Left, From, Below
|}
Bits 11 to 29 contain the data. Bit-field discrete data, Binary Coded Decimal (BCD), and Binary Number Representation (BNR) are common ARINC 429 data formats. Data formats may also be mixed.
Bits 9 and 10 are Source/Destination Identifiers (SDI) and may indicate the intended receiver or, more frequently, indicate the transmitting subsystem.
Bits 1 to 8 contain a label (label words), expressed in octal (MSB 1 bit numbering), identifying the data type.
The image below exemplifies many of the concepts explained in the adjacent sections. In this image the Label (260) appears in red, the Data in blue-green and the Parity bit in navy blue.
Labels
Label guidelines are provided as part of the ARINC 429 specification, for various equipment types. Each aircraft will contain a number of different systems, such as flight management computers, inertial reference systems, air data computers, radar altimeters, radios, and GPS sensors. For each type of equipment, a set of standard parameters is defined, which is common across all manufacturers and models. For example, any air data computer will provide the barometric altitude of the aircraft as label 203. This allows some degree of interchangeability of parts, as all air data computers behave, for the most part, in the same way. There are only a limited number of labels, though, and so label 203 may have some completely different meaning if sent by a GPS sensor, for example. Very commonly needed aircraft parameters, however, use the same label regardless of source. Also, as with any specification, each manufacturer has slight differences from the formal specification, such as by providing extra data above and beyond the specification, leaving out some data recommended by the specification, or other various changes.
Protection from interference
Avionics systems must meet environmental requirements, usually stated as RTCA DO-160 environmental categories. ARINC 429 employs several physical, electrical, and protocol techniques to minimize electromagnetic interference with on-board radios and other equipment, for example via other transmission cables.
Its cabling is a shielded 78 Ω twisted-pair. ARINC signaling defines a 10 Vp differential between the Data A and Data B levels within the bipolar transmission (i.e. 5 V on Data A and -5 V on Data B would constitute a valid driving signal), and the specification defines acceptable voltage rise and fall times.
ARINC 429's data encoding uses a complementary differential bipolar return-to-zero (BPRZ) transmission waveform, further reducing EMI emissions from the cable itself.
Development tools
When developing and/or troubleshooting the ARINC 429 bus, examination of hardware signals can be very important to find problems. A protocol analyzer is useful to collect, analyze, decode and store signals.
| Technology | Aircraft components | null |
3239191 | https://en.wikipedia.org/wiki/Post-transcriptional%20modification | Post-transcriptional modification | Transcriptional modification or co-transcriptional modification is a set of biological processes common to most eukaryotic cells by which an RNA primary transcript is chemically altered following transcription from a gene to produce a mature, functional RNA molecule that can then leave the nucleus and perform any of a variety of different functions in the cell. There are many types of post-transcriptional modifications achieved through a diverse class of molecular mechanisms.
One example is the conversion of precursor messenger RNA transcripts into mature messenger RNA that is subsequently capable of being translated into protein. This process includes three major steps that significantly modify the chemical structure of the RNA molecule: the addition of a 5' cap, the addition of a 3' polyadenylated tail, and RNA splicing. Such processing is vital for the correct translation of eukaryotic genomes because the initial precursor mRNA produced by transcription often contains both exons (coding sequences) and introns (non-coding sequences); splicing removes the introns and links the exons directly, while the cap and tail facilitate the transport of the mRNA to a ribosome and protect it from molecular degradation.
Post-transcriptional modifications may also occur during the processing of other transcripts which ultimately become transfer RNA, ribosomal RNA, or any of the other types of RNA used by the cell.
mRNA processing
5' processing
Capping
Capping of the pre-mRNA involves the addition of 7-methylguanosine (m7G) to the 5' end. To achieve this, the terminal 5' phosphate requires removal, which is done with the aid of enzyme RNA triphosphatase. The enzyme guanosyl transferase then catalyses the reaction, which produces the diphosphate 5' end. The diphosphate 5' end then attacks the alpha phosphorus atom of a GTP molecule in order to add the guanine residue in a 5'5' triphosphate link. The enzyme (guanine-N7-)-methyltransferase ("cap MTase") transfers a methyl group from S-adenosyl methionine to the guanine ring. This type of cap, with just the (m7G) in position is called a cap 0 structure. The ribose of the adjacent nucleotide may also be methylated to give a cap 1. Methylation of nucleotides downstream of the RNA molecule produce cap 2, cap 3 structures and so on. In these cases the methyl groups are added to the 2' OH groups of the ribose sugar.
The cap protects the 5' end of the primary RNA transcript from attack by ribonucleases that have specificity to the 3'5' phosphodiester bonds.
3' processing
Cleavage and polyadenylation
The pre-mRNA processing at the 3' end of the RNA molecule involves cleavage of its 3' end and then the addition of about 250 adenine residues to form a poly(A) tail. The cleavage and adenylation reactions occur primarily if a polyadenylation signal sequence (5'- AAUAAA-3') is located near the 3' end of the pre-mRNA molecule, which is followed by another sequence, which is usually (5'-CA-3') and is the site of cleavage. A GU-rich sequence is also usually present further downstream on the pre-mRNA molecule. More recently, it has been demonstrated that alternate signal sequences such as UGUA upstream off the cleavage site can also direct cleavage and polyadenylation in the absence of the AAUAAA signal. These two signals are not mutually independent, and often coexist. After the synthesis of the sequence elements, several multi-subunit proteins are transferred to the RNA molecule. The transfer of these sequence specific binding proteins cleavage and polyadenylation specificity factor (CPSF), Cleavage Factor I (CF I) and cleavage stimulation factor (CStF) occurs from RNA Polymerase II. The three factors bind to the sequence elements. The AAUAAA signal is directly bound by CPSF. For UGUA dependent processing sites, binding of the multi protein complex is done by Cleavage Factor I (CF I). The resultant protein complex formed contains additional cleavage factors and the enzyme Polyadenylate Polymerase (PAP). This complex cleaves the RNA between the polyadenylation sequence and the GU-rich sequence at the cleavage site marked by the (5'-CA-3') sequences. Poly(A) polymerase then adds about 200 adenine units to the new 3' end of the RNA molecule using ATP as a precursor. As the poly(A) tail is synthesized, it binds multiple copies of poly(A)-binding protein, which protects the 3'end from ribonuclease digestion by enzymes including the CCR4-Not complex.
Intron splicing
RNA splicing is the process by which introns, regions of RNA that do not code for proteins, are removed from the pre-mRNA and the remaining exons connected to re-form a single continuous molecule. Exons are sections of mRNA which become "expressed" or translated into a protein. They are the coding portions of a mRNA molecule. Although most RNA splicing occurs after the complete synthesis and end-capping of the pre-mRNA, transcripts with many exons can be spliced co-transcriptionally. The splicing reaction is catalyzed by a large protein complex called the spliceosome assembled from proteins and small nuclear RNA molecules that recognize splice sites in the pre-mRNA sequence. Many pre-mRNAs, including those encoding antibodies, can be spliced in multiple ways to produce different mature mRNAs that encode different protein sequences. This process is known as alternative splicing, and allows production of a large variety of proteins from a limited amount of DNA.
Histone mRNA processing
Histones H2A, H2B, H3 and H4 form the core of a nucleosome and thus are called core histones. Processing of core histones is done differently because typical histone mRNA lacks several features of other eukaryotic mRNAs, such as poly(A) tail and introns. Thus, such mRNAs do not undergo splicing and their 3' processing is done independent of most cleavage and polyadenylation factors. Core histone mRNAs have a special stem-loop structure at 3-prime end that is recognized by a stem–loop binding protein and a downstream sequence, called histone downstream element (HDE) that recruits U7 snRNA. Cleavage and polyadenylation specificity factor 73 cuts mRNA between stem-loop and HDE
Histone variants, such as H2A.Z or H3.3, however, have introns and are processed as normal mRNAs including splicing and polyadenylation.
| Biology and health sciences | Molecular biology | Biology |
3239291 | https://en.wikipedia.org/wiki/Fracture%20toughness | Fracture toughness | In materials science, fracture toughness is the critical stress intensity factor of a sharp crack where propagation of the crack suddenly becomes rapid and unlimited. A component's thickness affects the constraint conditions at the tip of a crack with thin components having plane stress conditions and thick components having plane strain conditions. Plane strain conditions give the lowest fracture toughness value which is a material property. The critical value of stress intensity factor in mode I loading measured under plane strain conditions is known as the plane strain fracture toughness, denoted . When a test fails to meet the thickness and other test requirements that are in place to ensure plane strain conditions, the fracture toughness value produced is given the designation . Fracture toughness is a quantitative way of expressing a material's resistance to crack propagation and standard values for a given material are generally available.
Slow self-sustaining crack propagation known as stress corrosion cracking, can occur in a corrosive environment above the threshold and below . Small increments of crack extension can also occur during fatigue crack growth, which after repeated loading cycles, can gradually grow a crack until final failure occurs by exceeding the fracture toughness.
Material variation
Fracture toughness varies by approximately 4 orders of magnitude across materials. Metals hold the highest values of fracture toughness. Cracks cannot easily propagate in tough materials, making metals highly resistant to cracking under stress and gives their stress–strain curve a large zone of plastic flow. Ceramics have a lower fracture toughness but show an exceptional improvement in the stress fracture that is attributed to their 1.5 orders of magnitude strength increase, relative to metals. The fracture toughness of composites, made by combining engineering ceramics with engineering polymers, greatly exceeds the individual fracture toughness of the constituent materials.
Mechanisms
Intrinsic mechanisms
Intrinsic toughening mechanisms are processes which act ahead of the crack tip to increase the material's toughness. These will tend to be related to the structure and bonding of the base material, as well as microstructural features and additives to it. Examples of mechanisms include:
crack deflection by secondary phases,
crack bifurcation due to fine grain structure
changes in the crack path due to grain boundaries
Any alteration to the base material which increases its ductility can also be thought of as intrinsic toughening.
Grain boundaries
The presence of grains in a material can also affect its toughness by affecting the way cracks propagate. In front of a crack, a plastic zone can be present as the material yields. Beyond that region, the material remains elastic. The conditions for fracture are the most favorable at the boundary between this plastic and elastic zone, and thus cracks often initiate by the cleavage of a grain at that location.
At low temperatures, where the material can become completely brittle, such as in a body-centered cubic (BCC) metal, the plastic zone shrinks away, and only the elastic zone exists. In this state, the crack will propagate by successive cleavage of the grains. At these low temperatures, the yield strength is high, but the fracture strain and crack tip radius of curvature are low, leading to a low toughness.
At higher temperatures, the yield strength decreases, and leads to the formation of the plastic zone. Cleavage is likely to initiate at the elastic-plastic zone boundary, and then link back to the main crack tip. This is usually a mixture of cleavages of grains, and ductile fracture of grains known as fibrous linkages. The percentage of fibrous linkages increase as temperature increases until the linkup is entirely fibrous linkages. In this state, even though yield strength is lower, the presence of ductile fracture and a higher crack tip radius of curvature results in a higher toughness.
Inclusions
Inclusions in a material such as a second phase particles can act similar to brittle grains that can affect crack propagation. Fracture or decohesion at the inclusion can either be caused by the external applied stress or by the dislocations generated by the requirement of the inclusion to maintain contiguity with the matrix around it. Similar to grains, the fracture is most likely to occur at the plastic-elastic zone boundary. Then the crack can linkup back to the main crack. If the plastic zone is small or the density of the inclusions is small, the fracture is more likely to directly link up with the main crack tip. If the plastic zone is large, or the density of inclusions is high, additional inclusion fractures may occur within the plastic zone, and linkup occurs by progressing from the crack to the closest fracturing inclusion within the zone.
Transformation toughening
Transformation toughening is a phenomenon whereby a material undergoes one or more martensitic (displacive, diffusionless) phase transformations which result in an almost instantaneous change in volume of that material. This transformation is triggered by a change in the stress state of the material, such as an increase in tensile stress, and acts in opposition to the applied stress. Thus when the material is locally put under tension, for example at the tip of a growing crack, it can undergo a phase transformation which increases its volume, lowering the local tensile stress and hindering the crack's progression through the material. This mechanism is exploited to increase the toughness of ceramic materials, most notably in Yttria-stabilized zirconia for applications such as ceramic knives and thermal barrier coatings on jet engine turbine blades.
Extrinsic mechanisms
Extrinsic toughening mechanisms are processes which act behind the crack tip to resist its further opening. Examples include
fibre/lamella bridging, where these structures hold the two fracture surfaces together after the crack has propagated through the matrix,
crack wedging from the friction between two rough fracture surfaces, and
microcracking, where smaller cracks form in the material around the main crack, relieving the stress at the crack tip by effectively increasing the material's compliance.
Test methods
Fracture toughness tests are performed to quantify the resistance of a material to failure by cracking. Such tests result in either a single-valued measure of fracture toughness or in a resistance curve. Resistance curves are plots where fracture toughness parameters (K, J etc.) are plotted against parameters characterizing the propagation of crack. The resistance curve or the single-valued fracture toughness is obtained based on the mechanism and stability of fracture. Fracture toughness is a critical mechanical property for engineering applications. There are several types of test used to measure fracture toughness of materials, which generally utilise a notched specimen in one of various configurations. A widely utilized standardized test method is the Charpy impact test whereby a sample with a V-notch or a U-notch is subjected to impact from behind the notch. Also widely used are crack displacement tests such as three-point beam bending tests with thin cracks preset into test specimens before applying load.
Testing requirements
Choice of specimen
The ASTM standard E1820 for the measurement of fracture toughness recommends three coupon types for fracture toughness testing, the single-edge bending coupon [SE(B)], the compact tension coupon [C(T)] and the disk-shaped compact tension coupon [DC(T)].
Each specimen configuration is characterized by three dimensions, namely the crack length (a), the thickness (B) and the width (W). The values of these dimensions are determined by the demand of the particular test that is being performed on the specimen. The vast majority of the tests are carried out on either compact or three-point flexural test configuration. For the same characteristic dimensions, compact configuration takes a lesser amount of material compared to three-point flexural test.
Material orientation
Orientation of fracture is important because of the inherent non-isotropic nature of most engineering materials. Due to this, there may be planes of weakness within the material, and crack growth along this plane may be easier compared to other direction. Due to this importance ASTM has devised a standardized way of reporting the crack orientation with respect to forging axis. The letters L, T and S are used to denote the longitudinal, transverse and short transverse directions, where the longitudinal direction coincides with forging axis. The orientation is defined with two letters the first one being the direction of principal tensile stress and the second one is the direction of crack propagation. Generally speaking, the lower bound of the toughness of a material is obtained in the orientation where the crack grows in the direction of forging axis.
Pre-cracking
For accurate results, a sharp crack is required before testing. Machined notches and slots do not meet this criterion. The most effective way of introducing a sufficiently sharp crack is by applying cyclic loading to grow a fatigue crack from a slot. Fatigue cracks are initiated at the tip of the slot and allowed to extend until the crack length reaches its desired value.
The cyclic loading is controlled carefully so as to not affect the toughness of the material through strain-hardening. This is done by choosing cyclic loads that produce a far smaller plastic zone compared to plastic zone of the main fracture. For example, according to ASTM E399, the maximum stress intensity Kmax should be no larger than 0.6 during the initial stage and less than 0.8 when crack approaches its final size.
In certain cases grooves are machined into the sides of a fracture toughness specimen so that the thickness of the specimen is reduced to a minimum of 80% of the original thickness along the intended path of crack extensions. The reason is to maintain a straight crack front during R-curve test.
The four main standardized tests are described below with KIc and KR tests valid for linear-elastic fracture mechanics (LEFM) while J and JR tests valid for elastic-plastic fracture mechanics (EPFM).
Plane-strain fracture toughness testing
When performing a fracture toughness test, the most common test specimen configurations are the single edge notch bend (SENB or three-point bend), and the compact tension (CT) specimens. Testing has shown that plane-strain conditions generally prevail when:
where is the minimum necessary thickness, the fracture toughness of the material and is the material yield strength.
The test is performed by loading steadily at a rate such that KI increases from 0.55 to 2.75 (MPa)/s. During the test, the load and the crack mouth opening displacement (CMOD) is recorded and the test is continued till the maximum load is reached. The critical load PQ is calculated through from the load vs CMOD plot. A provisional toughness KQ is given as
.
The geometry factor is a dimensionless function of a/W and is given in polynomial form in the E 399 standard. The geometry factor for compact test geometry can be found here. This provisional toughness value is recognized as valid when the following requirements are met:
and
When a material of unknown fracture toughness is tested, a specimen of full material section thickness is tested or the specimen is sized based on a prediction of the fracture toughness. If the fracture toughness value resulting from the test does not satisfy the requirement of the above equation, the test must be repeated using a thicker specimen. In addition to this thickness calculation, test specifications have several other requirements that must be met (such as the size of the shear lips) before a test can be said to have resulted in a KIC value.
When a test fails to meet the thickness and other plain-strain requirements, the fracture toughness value produced is given the designation Kc. Sometimes, it is not possible to produce a specimen that meets the thickness requirement. For example, when a relatively thin plate with high toughness is being tested, it might not be possible to produce a thicker specimen with plane-strain conditions at the crack tip.
Determination of R-curve, K-R
The specimen showing stable crack growth shows an increasing trend in fracture toughness as the crack length increases (ductile crack extension). This plot of fracture toughness vs crack length is called the resistance (R)-curve. ASTM E561 outlines a procedure for determining toughness vs crack growth curves in materials. This standard does not have a constraint over the minimum thickness of the material and hence can be used for thin sheets however the requirements for LEFM must be fulfilled for the test to be valid. The criteria for LEFM essentially states that in-plane dimension has to be large compared to the plastic zone. There is a misconception about the effect of thickness on the shape of R curve. It is hinted that for the same material thicker section fails by plane strain fracture and shows a single-valued fracture toughness, the thinner section fails by plane stress fracture and shows the rising R-curve. However, the main factor that controls the slope of R curve is the fracture morphology not the thickness. In some material section thickness changes the fracture morphology from ductile tearing to cleavage from thin to thick section, in which case the thickness alone dictates the slope of R-curve. There are cases where even plane strain fracture ensues in rising R-curve due to "microvoid coalescence" being the mode of failure.
The most accurate way of evaluating K-R curve is taking presence of plasticity into account depending on the relative size of the plastic zone. For the case of negligible plasticity, the load vs displacement curve is obtained from the test and on each point the compliance is found. The compliance is reciprocal of the slope of the curve that will be followed if the specimen is unloaded at a certain point, which can be given as the ratio of displacement to load for LEFM. The compliance is used to determine the instantaneous crack length through the relationship given in the ASTM standard.
The stress intensity should be corrected by calculating an effective crack length. ASTM standard suggests two alternative approaches. The first method is named Irwin's plastic zone correction. Irwin's approach describes the effective crack length to be
Irwin's approach leads to an iterative solution as K itself is a function of crack length.
The other method, namely the secant method, uses the compliance-crack length equation given by ASTM standard to calculate effective crack length from an effective compliance. Compliance at any point in Load vs displacement curve is essentially the reciprocal of the slope of the curve that ensues if the specimen is unloaded at that point. Now the unloading curve returns to the origin for linear elastic material but not for elastic plastic material as there is a permanent deformation. The effective compliance at a point for the elastic plastic case is taken as the slope of the line joining the point and origin (i.e the compliance if the material was an elastic one). This effective compliance is used to get an effective crack growth and the rest of the calculation follows the equation
The choice of plasticity correction is factored on the size of plastic zone. ASTM standard covering resistance curve suggests using Irwin's method is acceptable for small plastic zone and recommends using Secant method when crack-tip plasticity is more prominent. Also since the ASTM E 561 standard does not contain requirements on the specimen size or maximum allowable crack extension, thus the size independence of the resistance curve is not guaranteed. Few studies show that the size dependence is less detected in the experimental data for the Secant method.
Determination of JIC
Strain energy release rate per unit fracture surface area is calculated by J-integral method which is a contour path integral around the crack tip where the path begins and ends on either crack surfaces. J-toughness value signifies the resistance of the material in terms of amount of stress energy required for a crack to grow. JIC toughness value is measured for elastic-plastic materials. Now the single-valued JIC is determined as the toughness near the onset of the ductile crack extension (effect of strain hardening is not important). The test is performed with multiple specimen loading each of the specimen to various levels and unloading. This gives the crack mouth opening compliance which is to be used to get crack length with the help of relationships given in ASTM standard E 1820, which covers the J-integral testing. Another way of measuring crack growth is to mark the specimen with heat tinting or fatigue cracking. The specimen is eventually broken apart and the crack extension is measured with the help of the marks.
The test thus performed yields several load vs crack mouth opening displacement (CMOD) curves, which are used to calculate J as following:-
The linear elastic J is calculated using and K is determined from where is the net thickness for side-grooved specimen and equal to B for not side-grooved specimen.
The elastic plastic J is calculated using
Where
= 2 for SENB specimen
is initial ligament length given by the difference between width and initial crack length
is the plastic area under the load-displacement curve.
Specialized data reduction technique is used to get a provisional . The value is accepted if the following criterion is met:
Determination of tear resistance (Kahn tear test)
The tear test (e.g. Kahn tear test) provides a semi-quantitative measure of toughness in terms of tear resistance. This type of test requires a smaller specimen, and can, therefore, be used for a wider range of product forms. The tear test can also be used for very ductile aluminium alloys (e.g. 1100, 3003), where linear elastic fracture mechanics do not apply.
Standard test methods
A number of organizations publish standards related to fracture toughness measurements, namely ASTM, BSI, ISO, JSME.
ASTM C1161 Test Method for Flexural Strength of Advanced Ceramics at Ambient Temperature
ASTM C1421 Standard Test Methods for Determination of Fracture Toughness of Advanced Ceramics at Ambient Temperature
ASTM E399 Test Method for Plane-strain Fracture Toughness of Metallic Materials
ASTM E740 Practice for Fracture Testing with Surface-Crack Tension Specimens
ASTM E1820 Standard Test Method for Measurement of Fracture Toughness
ASTM E1823 Terminology Relating to Fatigue and Fracture Testing
ISO 12135 Metallic materials — Unified method of test for the determination of quasistatic fracture toughness
ISO 28079:2009, the Palmqvist method, used to determine the fracture toughness for cemented carbides
Crack deflection toughening
Many ceramics with polycrystalline structures develop large cracks that propagate along the boundaries between grains, rather than through the individual crystals themselves since the toughness of the grain boundaries is much lower than that of the crystals. The orientation of the grain boundary facets and residual stress cause the crack to advance in a complex, tortuous manner that is difficult to analyze. Simply calculating the additional surface energy associated with the increased grain boundary surface area due to this tortuosity is not accurate, as some of the energy to create the crack surface comes from the residual stress.
Faber–Evans model
A mechanics of materials model, introduced by Katherine Faber and Anthony G. Evans, has been developed to predict the increase in fracture toughness in ceramics due to crack deflection around second-phase particles that are prone to microcracking in a matrix. The model takes into account the particle morphology, aspect ratio, spacing, and volume fraction of the second phase, as well as the reduction in local stress intensity at the crack tip when the crack is deflected or the crack plane bows. The actual crack tortuosity is obtained through imaging techniques, allowing the deflection and bowing angles to be directly input into the model.
The resulting increase in fracture toughness is then compared to that of a flat crack through the plain matrix. The magnitude of the toughening is determined by the mismatch strain caused by thermal contraction incompatibility and the microfracture resistance of the particle/matrix interface. This toughening becomes noticeable when there is a narrow size distribution of particles that are appropriately sized. Researchers typically accept the findings of Faber's analysis, which suggest that deflection effects in materials with roughly equiaxial grains may increase the fracture toughness by about twice the grain boundary value.
| Physical sciences | Solid mechanics | Physics |
9157119 | https://en.wikipedia.org/wiki/Feasible%20region | Feasible region | In mathematical optimization and computer science, a feasible region, feasible set, or solution space is the set of all possible points (sets of values of the choice variables) of an optimization problem that satisfy the problem's constraints, potentially including inequalities, equalities, and integer constraints. This is the initial set of candidate solutions to the problem, before the set of candidates has been narrowed down.
For example, consider the problem of minimizing the function with respect to the variables and subject to and Here the feasible set is the set of pairs (x, y) in which the value of x is at least 1 and at most 10 and the value of y is at least 5 and at most 12. The feasible set of the problem is separate from the objective function, which states the criterion to be optimized and which in the above example is
In many problems, the feasible set reflects a constraint that one or more variables must be non-negative. In pure integer programming problems, the feasible set is the set of integers (or some subset thereof). In linear programming problems, the feasible set is a convex polytope: a region in multidimensional space whose boundaries are formed by hyperplanes and whose corners are vertices.
Constraint satisfaction is the process of finding a point in the feasible region.
Convex feasible set
A convex feasible set is one in which a line segment connecting any two feasible points goes through only other feasible points, and not through any points outside the feasible set. Convex feasible sets arise in many types of problems, including linear programming problems, and they are of particular interest because, if the problem has a convex objective function that is to be minimized, it will generally be easier to solve in the presence of a convex feasible set and any local optimum will also be a global optimum.
No feasible set
If the constraints of an optimization problem are mutually contradictory, there are no points that satisfy all the constraints and thus the feasible region is the empty set. In this case the problem has no solution and is said to be infeasible.
Bounded and unbounded feasible sets
Feasible sets may be bounded or unbounded. For example, the feasible set defined by the constraint set {x ≥ 0, y ≥ 0} is unbounded because in some directions there is no limit on how far one can go and still be in the feasible region. In contrast, the feasible set formed by the constraint set {x ≥ 0, y ≥ 0, x + 2y ≤ 4} is bounded because the extent of movement in any direction is limited by the constraints.
In linear programming problems with n variables, a necessary but insufficient condition for the feasible set to be bounded is that the number of constraints be at least n + 1 (as illustrated by the above example).
If the feasible set is unbounded, there may or may not be an optimum, depending on the specifics of the objective function. For example, if the feasible region is defined by the constraint set {x ≥ 0, y ≥ 0}, then the problem of maximizing x + y has no optimum since any candidate solution can be improved upon by increasing x or y; yet if the problem is to minimize x + y, then there is an optimum (specifically at (x, y) = (0, 0)).
Candidate solution
In optimization and other branches of mathematics, and in search algorithms (a topic in computer science), a candidate solution is a member of the set of possible solutions in the feasible region of a given problem. A candidate solution does not have to be a likely or reasonable solution to the problem—it is simply in the set that satisfies all constraints; that is, it is in the set of feasible solutions. Algorithms for solving various types of optimization problems often narrow the set of candidate solutions down to a subset of the feasible solutions, whose points remain as candidate solutions while the other feasible solutions are henceforth excluded as candidates.
The space of all candidate solutions, before any feasible points have been excluded, is called the feasible region, feasible set, search space, or solution space. This is the set of all possible solutions that satisfy the problem's constraints. Constraint satisfaction is the process of finding a point in the feasible set.
Genetic algorithm
In the case of the genetic algorithm, the candidate solutions are the individuals in the population being evolved by the algorithm.
Calculus
In calculus, an optimal solution is sought using the first derivative test: the first derivative of the function being optimized is equated to zero, and any values of the choice variable(s) that satisfy this equation are viewed as candidate solutions (while those that do not are ruled out as candidates). There are several ways in which a candidate solution might not be an actual solution. First, it might give a minimum when a maximum is being sought (or vice versa), and second, it might give neither a minimum nor a maximum but rather a saddle point or an inflection point, at which a temporary pause in the local rise or fall of the function occurs. Such candidate solutions may be able to be ruled out by use of the second derivative test, the satisfaction of which is sufficient for the candidate solution to be at least locally optimal. Third, a candidate solution may be a local optimum but not a global optimum.
In taking antiderivatives of monomials of the form the candidate solution using Cavalieri's quadrature formula would be This candidate solution is in fact correct except when
Linear programming
In the simplex method for solving linear programming problems, a vertex of the feasible polytope is selected as the initial candidate solution and is tested for optimality; if it is rejected as the optimum, an adjacent vertex is considered as the next candidate solution. This process is continued until a candidate solution is found to be the optimum.
| Mathematics | Optimization | null |
2349775 | https://en.wikipedia.org/wiki/Pinene | Pinene | Pinene is a collection of unsaturated bicyclic monoterpenes. Two geometric isomers of pinene are found in nature, α-pinene and β-pinene. Both are chiral. As the name suggests, pinenes are found in pines. Specifically, pinene is the major component of the liquid extracts of conifers. Pinenes are also found in many non-coniferous plants such as camphorweed (Heterotheca) and big sagebrush (Artemisia tridentata).
Isomers
Biosynthesis
α-Pinene and β-pinene are both produced from geranyl pyrophosphate, via cyclisation of linaloyl pyrophosphate followed by loss of a proton from the carbocation equivalent. Researchers at the Georgia Institute of Technology and the Joint BioEnergy Institute have been able to synthetically produce pinene with a bacterium.
Plants
Alpha-pinene is the most widely encountered terpenoid in nature and is highly repellent to insects.
Alpha-pinene appears in conifers and numerous other plants. Pinene is a major component of the essential oils of Sideritis spp. (ironwort) and
Salvia spp. (sage). Cannabis also contains alpha-pinene and beta-pinene. Resin from Pistacia terebinthus (commonly known as terebinth or turpentine tree) is rich in pinene. Pine nuts produced by pine trees contain pinene.
Makrut lime fruit peel contains an essential oil comparable to lime fruit peel oil; its main components are limonene and β-pinene.
The racemic mixture of the two forms of pinene is found in some oils like eucalyptus oil.
Reactions
α-Pinene
Selective oxidation of α-pinene occurs at the allylic position to give verbenone, along with pinene oxide, as well as verbenol and its hydroperoxide.
α-Pinene can be converted to camphor by way of isobornyl acetate.
Hydrogenation of pinene gives pinane, precursor to a useful pinanehydroperoxide.
The hydroboration of α-pinene has been extensively examined. With borane-dimethylsulfide, two equivalents of α-pinene react to give (diisopinocampheyl)borane. Reaction with 9-BBN gives the reagent called alpine borane. This sterically crowded chiral trialkylborane can stereoselectively reduce aldehydes in what is known as the Midland Alpine borane reduction.
β-Pinene
β-Pinene can be converted to α-pinene in the presence of strong bases, or pyrolysed to produce myrcene at 400 °C.
Use
Pinenes, especially α, are the primary constituents of turpentine, a nature-derived solvent and fuel.
The use of pinene as a biofuel in spark ignition engines has been explored. Pinene dimers have been shown to have heating values comparable to the jet fuel JP-10.
| Physical sciences | Terpenes and terpenoids | Chemistry |
2350490 | https://en.wikipedia.org/wiki/Road%20roller | Road roller | A road roller (sometimes called a roller-compactor, or just roller) is a compactor-type engineering vehicle used to compact soil, gravel, concrete, or asphalt in the construction of roads and foundations. Similar rollers are used also at landfills or in agriculture.
Road rollers are frequently referred to as steamrollers, regardless of their method of propulsion.
History
The first road rollers were horse-drawn, and were probably borrowed farm implements (see Roller).
Since the effectiveness of a roller depends to a large extent on its weight, self-powered vehicles replaced horse-drawn rollers from the mid-19th century. The first such vehicles were steam rollers. Single-cylinder steam rollers were generally used for base compaction and run with high engine revs with low gearing to promote bounce and vibration from the crankshaft through to the rolls in much the same way as a vibrating roller. The double cylinder or compound steam rollers became popular from around 1910 onwards and were used mainly for the rolling of hot-laid surfaces due to their smoother running engines, but both cylinder types are capable of rolling the finished surface. Steam rollers were often dedicated to a task by their gearing as the slower engines were for base compaction whereas the higher geared models were often referred to as "chip chasers" which followed the hot tar and chip laying machines. Some road companies in the US used steamrollers through the 1950s. In the UK some remained in service until the early 1970s.
As internal combustion engines improved during the 20th century, kerosene-, gasoline- (petrol), and diesel-powered rollers gradually replaced their steam-powered counterparts. The first internal-combustion powered road rollers were similar to the steam rollers they replaced. They used similar mechanisms to transmit power from the engine to the wheels, typically large, exposed spur gears. Some users disliked them in their infancy, as the engines of the era were typically difficult to start, particularly the kerosene-powered ones.
Virtually all road rollers in use today use diesel power.
Uses on a road
Road rollers use the weight of the vehicle to compress the surface being rolled (static) or use mechanical advantage (vibrating). Initial compaction of the substrate on a road project is done using a padfoot or "sheep's foot" drum roller, which achieves higher compaction density due to the pads having less surface area. On large freeways, a four-wheel compactor with padfoot drum and a blade, such as a Caterpillar 815/825 series machine, would be used due to its high weight, speed, and the powerful pushing force to spread bulk material. On regional roads, a smaller single padfoot drum machine may be used.
The next machine is usually a single smooth drum compactor that compacts the high spots down until the soil is smooth. This is usually done in combination with a motor grader to obtain a level surface. Sometimes at this stage a pneumatic tyre roller is used. These rollers feature two rows (front and back) of pneumatic tyres that overlap, and the flexibility of the tyres provides a kneading action that seals the surface and with some vertical movement of the wheels, enables the roller to operate effectively on uneven ground. Once the soil base is flat the pad drum compactor is no longer used on the road surface.
The next course (road base) is compacted using a smooth single drum, smooth tandem roller, or pneumatic tyre roller in combination with a grader and a water truck to achieve the desired flat surface with the correct moisture content for optimum compaction. Once the road base is compacted, the smooth single drum compactor is no longer used on the road surface (there is an exception if the single drum has special flat-wide-base tyres on the machine).
The final wear course of asphalt concrete (known as asphalt or blacktop in North America, or macadam in England) is laid using a paver and compacted using a tandem smooth drum roller, a three-point roller or a pneumatic tyre roller. Three point rollers on asphalt were once common and are still used, but tandem vibrating rollers are the usual choice now. The pneumatic tyre roller's kneading action is the final roller to seal the surface.
Rollers are also used in landfill compaction. Such compactors typically have padfoot drums, and do not achieve a smooth surface. The pads aid in compression, due to the smaller area contacting the ground.
Configurations
The roller can be a simple drum with a handle that is operated by one person and weighs or as large as a ride-on road roller weighing and costing more than US$150,000. A landfill unit may weigh .
Roller types
Pedestrian-operated
Rammer (bounce up and down)
Walk-behind plate compactor/light
Trench roller (manual unit or radio-frequency remote control)
Walk-behind roller/light (single drum)
Walk-behind roller/heavy (double drum)
Ride-on smooth finish
Tandem drum (static)
Tandem drum (vibrating)
Single drum roller (smooth)
Pneumatic-tyred Roller, called rubber tyre or multi-wheel
Combination roller (single row of tyres and a steel drum)
Three point roller (steam rollers are usually three-point)
Ride-on soil/landfill compactor with pads/feet/spikes
Single drum roller (soil)
4-wheel (soil/landfill)
3-point (soil/landfill)
Tandem drum (soil/landfill)
Other
Tractor-mounted and tractor-powered (conversion – see gallery picture below)
Drawn rollers or towed rollers (once common, now rare)
Impact compactor (uses a square or polygon drum to strike the ground hard for proof rolling or deep lift compacting)
Drum roller with rubber coated drum for asphalt compaction
Log skidder converted to compactor for landfill
Wheel loader converted to compactor for landfill
Drum types
Drums are available in widths ranging from .
Tyre roller types
Tyre rollers are available in widths ranging up to , with between 7 and 11 wheels (e.g. 3 wheels at front, 4 at back): 7 and 8 wheel types are normally used in Europe and Africa; 9 and 11 in America; and any type in Asia. Very heavy tyre rollers are used to compact soil.
Variations and features
On some machines, the drums may be filled with water on site to achieve the desired weight. When empty, the lighter machine is easier and cheaper to transport between work sites. On pneumatic tyre rollers the body may be ballasted with water or sand, or for extra compaction wet sand is used. Modern tyre rollers may be filled with steel ballast, which gives a more even balance for better compaction.
Additional compaction may be achieved by vibrating the roller drums, allowing a small, light machine to perform as well as a much heavier one. Vibration is typically produced by a free-spinning hydrostatic motor inside the drum to whose shaft an eccentric weight has been attached. Some rollers have a second weight that can be rotated relative to the main weight, to adjust the vibration amplitude and thus the compacting force.
Water lubrication may be provided to the drum surface from on-board "sprinkler tanks" to prevent hot asphalt sticking to the drum.
Hydraulic transmissions permit greater design flexibility. While early examples used direct mechanical drives, hydraulics reduce the number of moving parts exposed to contamination and allows the drum to be driven, providing extra traction on inclines.
Human-propelled rollers may only have a single roller drum.
Self-propelled rollers may have two drums, mounted one in front of the other (format known as "duplex"), or three rolls, or just one, with the back rollers replaced with treaded pneumatic tyres for increased traction.
Gallery
Manufacturers
ABG (Germany) — SD/TD (purchased by Ingersoll Rand and now part of Volvo CE)
Albaret (France) — PT/TD (now part of Caterpillar)
Ammann Group (Switzerland) —
Aveling-Barford (England) — TD/PT/3P
BOMAG (Germany) — SD/TD/PT (BOMAG/HYPAC in the US market)
Case CE (US) — SD (brands the Ammann/Sta machines as Case in the US)
Caterpillar Inc. (US) — SD/TD/PT (has the former lines of RAYGO, BROS and Bitelli)
Dynapac (Sweden) — SD/TD/PT/3P
Galion (US) —
Hamm AG (Germany) — SD/TD/PT/3P (now part of the Wirtgen Group)
HEPCO (Iran) — SD/TD/PT/3P
Hitachi (Japan) — SD/TD/PT/3P
Huber Company, (US) —
Hyster (US) — SD/TD/PT (part of HYPAC and Bomag USA)
Ingersoll Rand (US) — SD/TD/PT (now owned by Volvo)
Kamani Engineering Corporation (India) — tractor-mounted (now part of the RPG Group; production ended –1980)
LiuGong, HQ at Liuzhou, China —
Marshall (England) — TD
Sakai Heavy Industries (Japan) — SD/TD/PT/3P
Sany (China) — SD/TD/PT
World Equipment(China) — SD/TD/PT
Tampo (US) — SD/TD
Vibromax (Germany) — SD/TD/PT (purchased by JCB, now branded JCB)
Wacker Neuson (Germany) —
KEY:
SD = Single drum
TD = Tandem drum
PT = Pneumatic tyre — Rubber tyre or multi-tyre are also common
3P = 3-point rollers — These are very similar to the old steam roller design
| Technology | Specific-purpose transportation | null |
2350733 | https://en.wikipedia.org/wiki/Coma%20Cluster | Coma Cluster | The Coma Cluster (Abell 1656) is a large cluster of galaxies that contains over 1,000 identified galaxies. Along with the Leo Cluster (Abell 1367), it is one of the two major clusters comprising the Coma Supercluster. It is located in and takes its name from the constellation Coma Berenices.
The cluster's mean distance from Earth is 99 Mpc (321 million light years). Its ten brightest spiral galaxies have apparent magnitudes of 12–14 that are observable with amateur telescopes larger than 20 cm. The central region is dominated by two supergiant elliptical galaxies: NGC 4874 and NGC 4889. The cluster is within a few degrees of the north galactic pole on the sky. Most of the galaxies that inhabit the central portion of the Coma Cluster are ellipticals. Both dwarf and giant ellipticals are found in abundance in the Coma Cluster.
Cluster members
As is usual for clusters of this richness, the galaxies are overwhelmingly elliptical and S0 galaxies, with only a few spirals of younger age, and many of them probably near the outskirts of the cluster.
The full extent of the cluster was not understood until it was more thoroughly studied in the 1950s by astronomers at Mount Palomar Observatory, although many of the individual galaxies in the cluster had been identified previously.
Dark matter
The Coma Cluster is one of the first places where observed gravitational anomalies were considered to be indicative of unobserved mass. In 1933 Fritz Zwicky showed that the galaxies of the Coma Cluster were moving too fast for the cluster to be bound together by the visible matter of its galaxies. Though the idea of dark matter would not be accepted for another fifty years, Zwicky wrote that the galaxies must be held together by "dunkle Materie" (dark matter).
About 90% of the mass of the Coma cluster is believed to be in the form of dark matter. The distribution of dark matter throughout the cluster, however, is poorly constrained.
X-ray source
An extended X-ray source centered at 1300+28 in the direction of the Coma cluster of galaxies was reported before August 1966. This X-ray observation was performed by balloon, but the source was not detected in the sounding rocket flight launched by the X-ray astronomy group at the Naval Research Laboratory on November 25, 1964. A strong X-ray source was observed by the X-ray observatory satellite Uhuru close to the center of the Coma cluster and this source was suggested to be designated Coma X-1.
The Coma cluster contains about 800 galaxies within a 100 x 100 arc-min area of the celestial sphere. The source near the center at RA (1950) 12h56m ± 2m Dec 28°6' ± 12' has a luminosity Lx = 2.6 x 1044 ergs/s. As the source is extended, with a size of about 45', this argues against the possibility that a single galaxy is responsible for the emission. The Uhuru observations indicated a source strength of no greater than ~10−3 photons cm−2s−1keV−1 at 25 keV, which disagrees with the earlier observations claiming a source strength of ~10−2 photons cm−2s−1keV−1 at 25 keV, and a size of 5°.
Gallery
| Physical sciences | Notable galaxy clusters | Astronomy |
2350760 | https://en.wikipedia.org/wiki/Black%20mamba | Black mamba | The black mamba (Dendroaspis polylepis) is a species of highly venomous snake belonging to the family Elapidae. It is native to parts of sub-Saharan Africa. First formally described by Albert Günther in 1864, it is the second-longest venomous snake after the king cobra; mature specimens generally exceed and commonly grow to . Specimens of have been reported. It varies in colour from grey to dark brown. Juvenile black mambas tend to be paler than adults and darken with age. Despite the common name, the black mamba is not black; the colour name describes rather the inside of its mouth, which it displays when feeling threatened.
The species is both terrestrial (ground-living) and arboreal (tree-living); it inhabits savannah, woodland, rocky slopes and in some regions, dense forest. It is diurnal and is known to prey on birds and small mammals. Over suitable surfaces, it can move at speeds up to for short distances. Adult black mambas have few natural predators.
In a threat display, the black mamba usually opens its inky-black mouth, spreads its narrow neck-flap and sometimes hisses. It is capable of striking at considerable range and may deliver a series of bites in rapid succession. Its venom is primarily composed of neurotoxins that often induce symptoms within ten minutes, and is frequently fatal unless antivenom is administered. Despite its reputation as a formidable and highly aggressive species, the black mamba attacks humans only if it is threatened or cornered. It is rated as least concern on the International Union for Conservation of Nature (IUCN)'s Red List of Threatened Species.
Taxonomy
The first formal description of the black mamba was made in 1864 by German-born British zoologist Albert Günther. A single specimen was one of many species of snake collected by John Kirk, a naturalist who accompanied David Livingstone on the 1858–1864 Second Zambesi expedition. This specimen is the holotype and is housed in the Natural History Museum, London. The generic name of the species is derived from the Ancient Greek words dendron (δένδρον), "tree", and aspis (ἀσπίς) "asp", and the specific epithet polylepis is derived from the Ancient Greek poly (πολύ) meaning "many" and lepis (λεπίς) meaning "scale". The term "mamba" is derived from the Zulu word "imamba". In Tanzania, a local Ngindo name is ndemalunyayo ("grass-cutter") because it supposedly clips grass.
In 1873, German naturalist Wilhelm Peters described Dendraspis Antinorii from a specimen in the museum of Genoa that had been collected by Italian explorer Orazio Antinori in what is now northern Eritrea. This was subsequently regarded as a subspecies and is no longer held to be distinct. In 1896, Belgian-British zoologist George Albert Boulenger combined the species Dendroaspis polylepis as a whole with the eastern green mamba (Dendroaspis angusticeps), a lumping diagnosis that remained in force until 1946 when South African herpetologist Vivian FitzSimons again split them into separate species. A 2016 genetic analysis showed the black and eastern green mambas are each other's closest relatives, and are more distantly related to Jameson's mamba (Dendroaspis jamesoni), as shown in the cladogram below.
Description
The black mamba is a long, slender, cylindrical snake. It has a coffin-shaped head with a somewhat pronounced brow ridge and a medium-sized eyes. The adult snake's length typically ranges from but specimens have grown to lengths of . It is the longest species of venomous snake in Africa and the second-longest venomous snake species overall, exceeded in length only by the king cobra. The black mamba is a proteroglyphous (front-fanged) snake, with fangs up to in length, located at the front of the maxilla. The tail of the species is long and thin, the caudal vertebrae making up 17–25% of its body length. The body mass of black mambas has been reported to be about , although a study of seven black mambas found an average weight of , ranging from for a specimen of total length to for a specimen of total length.
Specimens vary considerably in colour, including olive, yellowish-brown, khaki and gunmetal but are rarely black. The scales of some individuals may have a purplish sheen. Individuals occasionally display dark mottling towards the posterior, which may appear in the form of diagonal crossbands. Black mambas have greyish-white underbellies. The common name is derived from the appearance of the inside of the mouth, dark bluish-grey to nearly black. Mamba eyes range between greyish-brown and shades of black; the pupil is surrounded by a silvery-white or yellow colour. Juvenile snakes are lighter in colour than adults; these are typically grey or olive green and darken as they age.
Scalation
The number and pattern of scales on a snake's body are a key element of identification to species level. The black mamba has between 23 and 25 rows of dorsal scales at midbody, 248 to 281 ventral scales, 109 to 132 divided subcaudal scales, and a divided anal scale. Its mouth is lined with 7–8 supralabial scales above, with the fourth and sometimes also the third one located under the eye, and 10-14 sublabial scales below. Its eyes have 3 or occasionally 4 preocular and 2–5 postocular scales.
Distribution and habitat
The black mamba inhabits a wide range in sub-Saharan Africa; its range includes Burkina Faso, Cameroon, Central African Republic, Democratic Republic of the Congo, South Sudan, Ethiopia, Eritrea, Somalia, Kenya, Uganda, Tanzania, Burundi, Rwanda, Mozambique, Eswatini, Malawi, Zambia, Zimbabwe, Botswana, South Africa, Namibia, and Angola. The black mamba's distribution in parts of West Africa has been disputed. In 1954, the black mamba was recorded in the Dakar region of Senegal. This observation, and a subsequent observation that identified a second specimen in the region in 1956, has not been confirmed and thus the snake's distribution in this area is inconclusive.
The species prefers moderately dry environments such as light woodland and scrub, rocky outcrops and semi-arid savanna. It also inhabits moist savanna and lowland forests. It is not commonly found at altitudes above , although its distribution does include locations at in Kenya and in Zambia. It is rated as a species of least concern on the International Union for Conservation of Nature (IUCN)'s Red List of endangered species, based on its huge range across sub-Saharan Africa and no documented decline.
Behaviour and ecology
The black mamba is both terrestrial and arboreal. On the ground, it moves with its head and neck raised, and typically uses termite mounds, abandoned burrows, rock crevices and tree cracks as shelter. Black mambas are diurnal; in South Africa, they are recorded to bask between 7 and 10 am and again from 2 to 4 pm. They may return daily to the same basking site.
Skittish and often unpredictable, the black mamba is agile and can move quickly. In the wild, black mambas seldom tolerate humans approaching more closely than about . When it perceives a threat, it retreats into brush or a hole. When confronted, it is likely to engage in a threat display, gaping to expose its black mouth and flicking its tongue. It also is likely to hiss and spread its neck into a hood similar to that of the cobras in the genus Naja.
During the threat display, any sudden movement by the intruder may provoke the snake into performing a series of rapid strikes, leading to severe envenomation. The size of the black mamba and its ability to raise its head a large distance from the ground enables it to launch as much as 40% of its body length upwards, so mamba bites to humans can occur on the upper body. The black mamba's reputation for being ready to attack is exaggerated; it is usually provoked by perceived threats such as the blocking of its movements and ability to retreat. The species' reputed speed has also been exaggerated; it cannot move more quickly than .
Reproduction and lifespan
The black mamba's breeding season spans from September to February, following the drop in temperature which occurs from April to June. Rival males compete by wrestling, attempting to subdue each other by intertwining their bodies and wrestling with their necks. Some observers have mistaken this for courtship. During mating, the male will slither over the dorsal side of the female while flicking his tongue. The female will signal her readiness to mate by lifting her tail and staying still. The male will then coil himself around the posterior end of the female and align his tail ventrolaterally with the female's. Intromission may last longer than two hours and the pair remain motionless apart from occasional spasms from the male.
The black mamba is oviparous; the female lays a clutch of 6–17 eggs. The eggs are elongated oval in shape, typically long and in diameter. When hatched, the young range from in length. They may grow quickly, reaching after their first year. Juvenile black mambas are very apprehensive and can be deadly like the adults. The black mamba is recorded to live up to 11 years and may live longer.
Feeding
The black mamba usually hunts from a permanent lair, to which it will regularly return if there is no disturbance. It mostly preys on small vertebrates such as birds, particularly nestlings and fledglings, and small mammals like rodents, bats, hyraxes and bushbabies. They generally prefer warm-blooded prey but will also consume other snakes. In the Transvaal area of South Africa, almost all recorded prey was rather small, largely consisting of rodents and similarly sized small or juvenile mammals as well as passerine birds, estimated to weigh only 1.9–7.8% of the mamba's body mass. Nonetheless, anecdotes have indicated that large black mambas may infrequently attack large prey such as the rock hyrax or dassie, and in some tribal languages, its name even means "dassie catcher". The black mamba does not typically hold onto its prey after biting; rather it releases its quarry and waits for it to succumb to paralysis and death before it is swallowed. The snake's potent digestive system has been recorded to fully digest prey in eight to ten hours.
Predation
Adult mambas have few natural predators aside from birds of prey. Brown snake eagles are verified predators of adult black mambas, of up to at least . Other eagles known to hunt or at least consume grown black mambas include tawny eagles and martial eagles. Young snakes have been recorded as prey of the Cape file snake. Mongooses, which have some resistance to mamba venom and are often quick enough to evade a bite, will sometimes harass or take a black mamba for prey, and may pursue them in trees. The similarly predatory honey badger also has some resistance to mamba venom. The mechanism in both mammals is thought to be that their muscular nicotinic acetylcholine receptors do not bind snake alpha-neurotoxins. Black mambas have also been found amongst the stomach contents of Nile crocodiles. Young mambas in the Serengeti are known to fall prey to southern ground hornbills, marsh owls and hooded vultures.
Venom
The black mamba is the most feared snake in Africa because of its size, aggression, venom toxicity and speed of onset of symptoms following envenomation, and is classified as a snake of medical importance by the World Health Organization. A survey in South Africa from 1957 to 1979 recorded 2,553 venomous snakebites, 75 of which were confirmed as being from black mambas. Of these 75 cases, 63 had symptoms of systemic envenomation and 21 died. Those bitten before 1962 received a polyvalent antivenom that had no effect on black mamba venom, and 15 of 35 people who received the antivenom died. A mamba-specific antivenom was introduced in 1962, followed by a fully polyvalent antivenom in 1971. Over this period, 5 of 38 people bitten by black mambas and given antivenom died. A census in rural Zimbabwe in 1991 and 1992 revealed 274 cases of snakebite, of which 5 died. Black mambas were confirmed in 15 cases, of which 2 died. The peak period for deaths is the species' breeding season from September to February, during which black mambas are most irritable. Bites are very rare outside Africa; snake handlers and enthusiasts are the usual victims.
Unlike many venomous snake species, black mamba venom does not contain protease enzymes. Its bites do not generally cause local swelling or necrosis, and the only initial symptom may be a tingling sensation in the area of the bite. The snake tends to bite repeatedly and let go, so there can be multiple puncture wounds. Its bite can deliver about 100–120 mg of venom on average; the maximum recorded dose is 400 mg. The murine median lethal dose (LD50) when administered intravenously has been calculated at 0.32 and 0.33 mg/kg. Bites were often fatal before antivenom was widely available.
The venom is predominantly neurotoxic, and symptoms often become apparent within 10 minutes. Early neurological signs that indicate severe envenomation include a metallic taste, drooping eyelids (ptosis) and gradual symptoms of bulbar palsy. Other neurological symptoms include miosis (constricted pupils), blurred or diminished vision, paresthesia (a tingling sensation on the skin), dysarthria (slurred speech), dysphagia (difficulty swallowing), dyspnea (shortness of breath), difficulty handling saliva, an absent gag reflex, fasciculations (muscle twitches), ataxia (impaired voluntary movement), vertigo, drowsiness and loss of consciousness, and respiratory paralysis. Other more general symptoms include nausea and vomiting, abdominal pain, diarrhea, sweating, salivation, goosebumps and red eyes. The bite of a black mamba can cause collapse in humans within 45 minutes. Without appropriate antivenom treatment, symptoms typically progress to respiratory failure, which leads to cardiovascular collapse and death. This typically occurs in 7 to 15 hours.
In 2015, the proteome (complete protein profile) of black mamba venom was assessed and published, revealing 41 distinct proteins and one nucleoside. The venom is composed of two main families of toxic agents, dendrotoxins (I and K) and (at a slightly lower proportion) three-finger toxins. Dendrotoxins are akin to kunitz-type protease inhibitors that interact with voltage-dependent potassium channels, stimulating acetylcholine and causing an excitatory effect, and are thought to cause symptoms such as sweating. Members of the three-finger family include alpha-neurotoxin, cardiotoxins, fasciculins and mambalgins. The most toxic components are the alpha-neurotoxins, which bind nicotinic acetylcholine receptors and hence block the action of acetylcholine at the postsynaptic membrane and cause neuromuscular blockade and hence paralysis. Fasciculins are anticholinesterase inhibitors that cause muscle fasciculation. The venom has little or no haemolytic, haemorrhagic or procoagulant activity. Mambalgins act as inhibitors for acid-sensing ion channels in the central and peripheral nervous system, causing a pain-inhibiting effect. There is research interest in their analgesic potential.
The composition of black mamba venom differs markedly from those of other mambas, all of which contain predominantly three-finger toxin agents. It is thought this may reflect the preferred prey items – small mammals for the mainly land-dwelling black mamba versus birds for the other predominantly arboreal mambas. Unlike many snake species, black mamba venom has little phospholipase A2 content.
Treatment
Standard first aid treatment for any suspected bite from a venomous snake is the application of a pressure bandage to the bite site, minimisation of movement of the victim and conveyance to a hospital or clinic as quickly as possible. The neurotoxic nature of black mamba venom means an arterial tourniquet may be of benefit. Tetanus toxoid is sometimes administered, though the main treatment is the administration of the appropriate antivenom. A polyvalent antivenom produced by the South African Institute for Medical Research is used to treat black mamba bites, and a new antivenom was being developed by the Universidad de Costa Rica's Instituto Clodomiro Picado.
Notable bite cases
Danie Pienaar, who was at various times from at least 2009 to 2017 head of South African National Parks Scientific Services and acting managing executive, survived the bite of a black mamba without antivenom in 1998. Despite the hospital physicians having declared it a "moderate" envenomation, Pienaar lapsed into a coma at one point and his prognosis was declared "poor". Upon arrival at the hospital, Pienaar was immediately intubated and placed on life support for 3 days. He was released from the hospital on the fifth day. Remaining calm after being bitten increased his chances of survival, as did the application of a tourniquet.
In March 2008, 28-year-old British trainee safari guide Nathan Layton was bitten by a black mamba that had been found near his classroom at the Southern African Wildlife College in Hoedspruit, Limpopo, South Africa. Layton was bitten by the snake on his index finger while it was being put into a jar and first aid-trained staff who examined him determined he could carry on with lectures. He thought the snake had only brushed his hand. Layton complained of blurred vision within an hour of being bitten, and collapsed and died shortly afterwards.
American professional photographer Mark Laita was bitten on the leg by a black mamba during a photo-shoot of a black mamba at a facility in Central America. Bleeding profusely, he did not seek medical attention, and except for intense pain and local swelling overnight, he was not affected. This led him to believe that either the snake gave him a "dry bite" (a bite without injecting venom) or the heavy bleeding pushed the venom out. Some commenters on the story suggested that it was a venomoid snake (in which the venom glands are surgically removed), but Laita responded that it was not. Only later did Laita find that he had captured the snake biting his leg in a photograph.
In 2016, Kenyan woman Cheposait Adomo was attacked by three black mambas, one of which bit her repeatedly on the leg, in West Pokot County, Kenya. People coming to her aid drove off the other snakes, hacking two with a machete. After an attempt at using traditional medicine, they placed her on a motorcycle and conveyed her 45 minutes to the nearest hospital, which had antivenom. She survived.
Prominent South African anti-Apartheid activist and Labour Court judge Anton Steenkamp died after being bitten by a black mamba while on leave in Zambia in May 2019. He was several hours away from medical help and died before antivenom could be administered.
In June 2020, Bulgarian veterinarian Georgi Elenski from Haskovo was bitten by a black mamba that was part of his personal collection of exotic animals. His initial condition was very serious, but he recovered after extensive treatment involving the administering of antivenom and respiratory support.
In January 2022, a former newspaper office worker and farmer from Zimbabwe, Peter Dube, died after getting bitten by a black mamba, due to the hospital he was taken to not having any antivenom to treat him.
In January 2023, a 17-year-old student from Zimbabwe died after being bitten by a black mamba. The snake had gone into a high school classroom while the students were outside.
| Biology and health sciences | Snakes | Animals |
2350918 | https://en.wikipedia.org/wiki/Thermal%20runaway | Thermal runaway | Thermal runaway describes a process that is accelerated by increased temperature, in turn releasing energy that further increases temperature. Thermal runaway occurs in situations where an increase in temperature changes the conditions in a way that causes a further increase in temperature, often leading to a destructive result. It is a kind of uncontrolled positive feedback.
In chemistry (and chemical engineering), thermal runaway is associated with strongly exothermic reactions that are accelerated by temperature rise. In electrical engineering, thermal runaway is typically associated with increased current flow and power dissipation. Thermal runaway can occur in civil engineering, notably when the heat released by large amounts of curing concrete is not controlled. In astrophysics, runaway nuclear fusion reactions in stars can lead to nova and several types of supernova explosions, and also occur as a less dramatic event in the normal evolution of solar-mass stars, the "helium flash".
Chemical engineering
Chemical reactions involving thermal runaway are also called thermal explosions in chemical engineering, or runaway reactions in organic chemistry. It is a process by which an exothermic reaction goes out of control: the reaction rate increases due to an increase in temperature, causing a further increase in temperature and hence a further rapid increase in the reaction rate. This has contributed to industrial chemical accidents, most notably the 1947 Texas City disaster from overheated ammonium nitrate in a ship's hold, and the 1976 explosion of zoalene, in a drier, at King's Lynn. Frank-Kamenetskii theory provides a simplified analytical model for thermal explosion. Chain branching is an additional positive feedback mechanism which may also cause temperature to skyrocket because of rapidly increasing reaction rate.
Chemical reactions are either endothermic or exothermic, as expressed by their change in enthalpy. Many reactions are highly exothermic, so many industrial-scale and oil refinery processes have some level of risk of thermal runaway. These include hydrocracking, hydrogenation, alkylation (SN2), oxidation, metalation and nucleophilic aromatic substitution. For example, oxidation of cyclohexane into cyclohexanol and cyclohexanone and ortho-xylene into phthalic anhydride have led to catastrophic explosions when reaction control failed.
Thermal runaway may result from unwanted exothermic side reaction(s) that begin at higher temperatures, following an initial accidental overheating of the reaction mixture. This scenario was behind the Seveso disaster, where thermal runaway heated a reaction to temperatures such that in addition to the intended 2,4,5-trichlorophenol, poisonous 2,3,7,8-tetrachlorodibenzo-p-dioxin was also produced, and was vented into the environment after the reactor's rupture disk burst.
Thermal runaway is most often caused by failure of the reactor vessel's cooling system. Failure of the mixer can result in localized heating, which initiates thermal runaway. Similarly, in flow reactors, localized insufficient mixing causes hotspots to form, wherein thermal runaway conditions occur, which causes violent blowouts of reactor contents and catalysts. Incorrect equipment component installation is also a common cause. Many chemical production facilities are designed with high-volume emergency venting, a measure to limit the extent of injury and property damage when such accidents occur.
At large scale, it is unsafe to "charge all reagents and mix", as is done in laboratory scale. This is because the amount of reaction scales with the cube of the size of the vessel (V ∝ r³), but the heat transfer area scales with the square of the size (A ∝ r²), so that the heat production-to-area ratio scales with the size (V/A ∝ r). Consequently, reactions that easily cool fast enough in the laboratory can dangerously self-heat at ton scale. In 2007, this kind of erroneous procedure caused an explosion of a -reactor used to metalate methylcyclopentadiene with metallic sodium, causing the loss of four lives and parts of the reactor being flung away. Thus, industrial scale reactions prone to thermal runaway are preferably controlled by the addition of one reagent at a rate corresponding to the available cooling capacity.
Some laboratory reactions must be run under extreme cooling, because they are very prone to hazardous thermal runaway. For example, in Swern oxidation, the formation of sulfonium chloride must be performed in a cooled system (−30 °C), because at room temperature the reaction undergoes explosive thermal runaway.
Microwave heating
Microwaves are used for heating of various materials in cooking and various industrial processes. The rate of heating of the material depends on the energy absorption, which depends on the dielectric constant of the material. The dependence of dielectric constant on temperature varies for different materials; some materials display significant increase with increasing temperature. This behavior, when the material gets exposed to microwaves, leads to selective local overheating, as the warmer areas are better able to accept further energy than the colder areas—potentially dangerous especially for thermal insulators, where the heat exchange between the hot spots and the rest of the material is slow. These materials are called thermal runaway materials. This phenomenon occurs in some ceramics.
Electrical engineering
Some electronic components develop lower resistances or lower triggering voltages (for nonlinear resistances) as their internal temperature increases. If circuit conditions cause markedly increased current flow in these situations, increased power dissipation may raise the temperature further by Joule heating. A vicious circle or positive feedback effect of thermal runaway can cause failure, sometimes in a spectacular fashion (e.g. electrical explosion or fire). To prevent these hazards, well-designed electronic systems typically incorporate current limiting protection, such as thermal fuses, circuit breakers, or PTC current limiters.
To handle larger currents, circuit designers may connect multiple lower-capacity devices (e.g. transistors, diodes, or MOVs) in parallel. This technique can work well, but is susceptible to a phenomenon called current hogging, in which the current is not shared equally across all devices. Typically, one device may have a slightly lower resistance, and thus draws more current, heating it more than its sibling devices, causing its resistance to drop further. The electrical load ends up funneling into a single device, which then rapidly fails. Thus, an array of devices may end up no more robust than its weakest component.
The current-hogging effect can be reduced by carefully matching the characteristics of each paralleled device, or by using other design techniques to balance the electrical load. However, maintaining load balance under extreme conditions may not be straightforward. Devices with an intrinsic positive temperature coefficient (PTC) of electrical resistance are less prone to current hogging, but thermal runaway can still occur because of poor heat sinking or other problems.
Many electronic circuits contain special provisions to prevent thermal runaway. This is most often seen in transistor biasing arrangements for high-power output stages. However, when equipment is used above its designed ambient temperature, thermal runaway can still occur in some cases. This occasionally causes equipment failures in hot environments, or when air cooling vents are blocked.
Semiconductors
Silicon shows a peculiar profile, in that its electrical resistance increases with temperature up to about 160 °C, then starts decreasing, and drops further when the melting point is reached. This can lead to thermal runaway phenomena within internal regions of the semiconductor junction; the resistance decreases in the regions which become heated above this threshold, allowing more current to flow through the overheated regions, in turn causing yet more heating in comparison with the surrounding regions, which leads to further temperature increase and resistance decrease. This leads to the phenomenon of current crowding and formation of current filaments (similar to current hogging, but within a single device), and is one of the underlying causes of many semiconductor junction failures.
Bipolar junction transistors (BJTs)
Leakage current increases significantly in bipolar transistors (especially germanium-based bipolar transistors) as they increase in temperature. Depending on the design of the circuit, this increase in leakage current can increase the current flowing through a transistor and thus the power dissipation, causing a further increase in collector-to-emitter leakage current. This is frequently seen in a push–pull stage of a class AB amplifier. If the pull-up and pull-down transistors are biased to have minimal crossover distortion at room temperature, and the biasing is not temperature-compensated, then as the temperature rises both transistors will be increasingly biased on, causing current and power to further increase, and eventually destroying one or both devices.
One rule of thumb to avoid thermal runaway is to keep the operating point of a BJT so that Vce ≤ 1/2 Vcc
Another practice is to mount a thermal feedback sensing transistor or other device on the heat sink, to control the crossover bias voltage. As the output transistors heat up, so does the thermal feedback transistor. This in turn causes the thermal feedback transistor to turn on at a slightly lower voltage, reducing the crossover bias voltage, and so reducing the heat dissipated by the output transistors.
If multiple BJT transistors are connected in parallel (which is typical in high current applications), a current hogging problem can occur. Special measures must be taken to control this characteristic vulnerability of BJTs.
In power transistors (which effectively consist of many small transistors in parallel), current hogging can occur between different parts of the transistor itself, with one part of the transistor becoming more hot than the others. This is called second breakdown, and can result in destruction of the transistor even when the average junction temperature seems to be at a safe level.
Power MOSFETs
Power MOSFETs typically increase their on-resistance with temperature. Under some circumstances, power dissipated in this resistance causes more heating of the junction, which further increases the junction temperature, in a positive feedback loop. As a consequence, power MOSFETs have stable and unstable regions of operation. However, the increase of on-resistance with temperature helps balance current across multiple MOSFETs connected in parallel, so current hogging does not occur. If a MOSFET transistor produces more heat than the heatsink can dissipate, then thermal runaway can still destroy the transistors. This problem can be alleviated to a degree by lowering the thermal resistance between the transistor die and the heatsink. | Physical sciences | Thermodynamics | Chemistry |
2351790 | https://en.wikipedia.org/wiki/Crocoite | Crocoite | Crocoite is a mineral consisting of lead chromate, PbCrO4, and crystallizing in the monoclinic crystal system. It is identical in composition with the artificial product chrome yellow used as a paint pigment.
Description
Crocoite is commonly found as large, well-developed prismatic adamantine crystals, although in many cases are poorly terminated. Crystals are of a bright hyacinth-red color, translucent, and have an adamantine to vitreous lustre. On exposure to UV light some of the translucency and brilliancy is lost. The streak is orange-yellow; Mohs hardness is 2.5–3; and the specific gravity is 6.0.
It was discovered at the Berezovskoe Au Deposit (Berezovsk Mines) near Ekaterinburg in the Urals in 1766; and named crocoise by F. S. Beudant in 1832, from the Greek κρόκος (krokos), saffron, in allusion to its color, a name first altered to crocoisite and afterwards to crocoite. In the type locality the crystals are found in gold-bearing quartz-veins traversing granite or gneiss and associated with crocoite are quartz, embreyite, phoenicochroite and vauquelinite. Phoenicochroite is a basic lead chromate, Pb2CrO5 with dark red crystals, and vauquelinite a lead and copper phosphate-chromate, Pb2CuCrO4PO4OH, with brown or green monoclinic crystals. Vauquelinite was named after Louis Nicolas Vauquelin, who in 1797 discovered (simultaneously with and independently of M. H. Klaproth) the element chromium in crocoite.
Abundant masses with exceptional examples of crocoite crystals have been found in the Extended Mine at Mount Dundas as well as the Adelaide, Red Lead, West Comet, Platt and a few other Mines at Dundas, Tasmania; they are usually found in long slender prisms, usually about 10–20 mm but rarely up to 100 mm (4 inches) in length, with a brilliant lustre and color. Crocoite is also the official Tasmanian mineral emblem. Other localities which have yielded good crystallized specimens are Congonhas do Campo near Ouro Preto in Brazil, Luzon in the Philippines, Mutare in Mashonaland, near Menzies in Western Australia, plus Brazil, Germany and South Africa.
The relative rarity of crocoite is connected with the specific conditions required for its formation: an oxidation zone of lead ore bed and presence of ultramafic rocks serving as the source of chromium (in chromite). Oxidation of Cr3+ into CrO42− (from chromite) and decomposition of galena (or other primary lead minerals) are required for crocoite formation. These conditions are relatively unusual.
As crocoite is composed of lead(II) chromate, it is toxic, containing both lead and hexavalent chromium.
Crocoite from Tasmania has been mined from the Dundas Extended Mine by Mike and Eleanor Phelan since the mid-1980s, but the mine's origins date back to 1892 when it was used as a prospecting tunnel for silver lead. As at April 2019, the mine is for sale (A$300,000) with the owners then continuing to operate the nearby Stichtite mine.
Gallery
Examples of crocoite
| Physical sciences | Minerals | Earth science |
2351950 | https://en.wikipedia.org/wiki/Catty | Catty | The catty or kati is a traditional Chinese unit of mass used across East and Southeast Asia, notably for weighing food and other groceries. Related units include the picul, equal to 100 catties, and the tael, which is of a catty. A stone is a former unit used in Hong Kong equal to 120 catties and a gwan () is 30 catties. Catty or kati is still used in Southeast Asia as a unit of measurement in some contexts especially by the significant Overseas Chinese populations across the region, particularly in Malaysia and Singapore.
The catty is traditionally equivalent to around pound avoirdupois, formalised as 604.78982 grams in Hong Kong, 604.5 grams historically in Vietnam, 604.79 grams in Malaysia and 604.8 grams in Singapore. In some countries, the weight has been rounded to 600 grams (Taiwan, Japan, Korea and Thailand). In mainland China, the catty (more commonly translated as jin within China) has been rounded to 500 grams and is referred to as the market catty ( ) in order to distinguish it from the kilogram, called the common catty ( ), and it is subdivided into 10 taels rather than the usual 16.
Etymology
The word catty comes from Malay kati, meaning 'the weight'. It has also been borrowed into English as caddy, meaning a container for storing tea.
Gallery
| Physical sciences | Chinese | Basics and measurement |
2352847 | https://en.wikipedia.org/wiki/Strategic%20dominance | Strategic dominance | In game theory, a strategy A dominates another strategy B if A will always produces a better result than B, regardless of how any other player plays no matter how that player's opponent or opponents play. Some very simple games (called straightforward games) can be solved using dominance.
Terminology
A player can compare two strategies, A and B, to determine which one is better. The result of the comparison is one of:
B strictly dominates (>) A: choosing B always gives a better outcome than choosing A, no matter what the other players do.
B weakly dominates (≥) A: choosing B always gives at least as good an outcome as choosing A, no matter what the other players do, and there is at least one set of opponents' actions for which B gives a better outcome than A. (Notice that if B strictly dominates A, then B weakly dominates A. Therefore, we can say "B dominates A" to mean "B weakly dominates A".)
B is weakly dominated by A: there is at least one set of opponents' actions for which B gives a worse outcome than A, while all other sets of opponents' actions give B the same payoff as A. (Strategy A weakly dominates B).
B is strictly dominated by A: choosing B always gives a worse outcome than choosing A, no matter what the other player(s) do. (Strategy A strictly dominates B).
Neither A nor B dominates the other: B and A are not equivalent, and B neither dominates, nor is dominated by, A. Choosing A is better in some cases, while choosing B is better in other cases, depending on exactly how the opponent chooses to play. For example, B is "throw rock" while A is "throw scissors" in Rock, Paper, Scissors.
This notion can be generalized beyond the comparison of two strategies.
Strategy B is strictly dominant if strategy B strictly dominates every other possible strategy.
Strategy B is weakly dominant if strategy B weakly dominates every other possible strategy.
Strategy B is strictly dominated if some other strategy exists that strictly dominates B.
Strategy B is weakly dominated if some other strategy exists that weakly dominates B.
Strategy: A complete contingent plan for a player in the game. A complete contingent plan is a full specification of a player's behavior, describing each action a player would take at every possible decision point. Because information sets represent points in a game where a player must make a decision, a player's strategy describes what that player will do at each information set.
Rationality: The assumption that each player acts in a way that is designed to bring about what he or she most prefers given probabilities of various outcomes; von Neumann and Morgenstern showed that if these preferences satisfy certain conditions, this is mathematically equivalent to maximizing a payoff. A straightforward example of maximizing payoff is that of monetary gain, but for the purpose of a game theory analysis, this payoff can take any desired outcome—cash reward, minimization of exertion or discomfort, or promoting justice can all be modeled as amassing an overall “utility” for the player. The assumption of rationality states that players will always act in the way that best satisfies their ordering from best to worst of various possible outcomes.
Common Knowledge: The assumption that each player has knowledge of the game, knows the rules and payoffs associated with each course of action, and realizes that every other player has this same level of understanding. This is the premise that allows a player to make a value judgment on the actions of another player, backed by the assumption of rationality, into consideration when selecting an action.
Dominance and Nash equilibria
If a strictly dominant strategy exists for one player in a game, that player will play that strategy in each of the game's Nash equilibria. If both players have a strictly dominant strategy, the game has only one unique Nash equilibrium, referred to as a "dominant strategy equilibrium". However, that Nash equilibrium is not necessarily "efficient", meaning that there may be non-equilibrium outcomes of the game that would be better for both players. The classic game used to illustrate this is the Prisoner's Dilemma.
Strictly dominated strategies cannot be a part of a Nash equilibrium, and as such, it is irrational for any player to play them. On the other hand, weakly dominated strategies may be part of Nash equilibria. For instance, consider the payoff matrix pictured at the right.
Strategy C weakly dominates strategy D. Consider playing C: If one's opponent plays C, one gets 1; if one's opponent plays D, one gets 0. Compare this to D, where one gets 0 regardless. Since in one case, one does better by playing C instead of D and never does worse, C weakly dominates D. Despite this, is a Nash equilibrium. Suppose both players choose D. Neither player will do any better by unilaterally deviating—if a player switches to playing C, they will still get 0. This satisfies the requirements of a Nash equilibrium. Suppose both players choose C. Neither player will do better by unilaterally deviating—if a player switches to playing D, they will get 0. This also satisfies the requirements of a Nash equilibrium.
Iterated elimination of strictly dominated strategies
The iterated elimination (or deletion, or removal) of dominated strategies (also denominated as IESDS, or IDSDS, or IRSDS) is one common technique for solving games that involves iteratively removing dominated strategies. In the first step, all dominated strategies are removed from the strategy space of each of the players, since no rational player would ever play these strategies. This results in a new, smaller game. Some strategies—that were not dominated before—may be dominated in the smaller game. The first step is repeated, creating a new even smaller game, and so on.
This process is valid since it is assumed that rationality among players is common knowledge, that is, each player knows that the rest of the players are rational, and each player knows that the rest of the players know that he knows that the rest of the players are rational, and so on ad infinitum (see Aumann, 1976).
| Mathematics | Game theory | null |
2352910 | https://en.wikipedia.org/wiki/Solar%20cell | Solar cell | A solar cell, also known as a photovoltaic cell (PV cell), is an electronic device that converts the energy of light directly into electricity by means of the photovoltaic effect. It is a form of photoelectric cell, a device whose electrical characteristics (such as current, voltage, or resistance) vary when it is exposed to light. Individual solar cell devices are often the electrical building blocks of photovoltaic modules, known colloquially as "solar panels". Almost all commercial PV cells consist of crystalline silicon, with a market share of 95%. Cadmium telluride thin-film solar cells account for the remainder. The common single-junction silicon solar cell can produce a maximum open-circuit voltage of approximately 0.5 to 0.6 volts.
Photovoltaic cells may operate under sunlight or artificial light. In addition to producing energy, they can be used as a photodetector (for example infrared detectors), detecting light or other electromagnetic radiation near the visible range, or measuring light intensity.
The operation of a PV cell requires three basic attributes:
The absorption of light, generating excitons (bound electron-hole pairs), unbound electron-hole pairs (via excitons), or plasmons.
The separation of charge carriers of opposite types.
The separate extraction of those carriers to an external circuit.
In contrast, a solar thermal collector supplies heat by absorbing sunlight, for the purpose of either direct heating or indirect electrical power generation from heat. A "photoelectrolytic cell" (photoelectrochemical cell), on the other hand, refers either to a type of photovoltaic cell (like that developed by Edmond Becquerel and modern dye-sensitized solar cells), or to a device that splits water directly into hydrogen and oxygen using only solar illumination.
Photovoltaic cells and solar collectors are the two means of producing solar power.
Applications
Assemblies of solar cells are used to make solar modules that generate electrical power from sunlight, as distinguished from a "solar thermal module" or "solar hot water panel". A solar array generates solar power using solar energy.
Vehicular applications
Application of solar cells as an alternative energy source for vehicular applications is a growing industry. Electric vehicles that operate off of solar energy and/or sunlight are commonly referred to as solar cars. These vehicles use solar panels to convert absorbed light into electrical energy that is then stored in batteries. There are multiple input factors that affect the output power of solar cells such as temperature, material properties, weather conditions, solar irradiance and more.
The first instance of photovoltaic cells within vehicular applications was around midway through the second half of the 1900's. In an effort to increase publicity and awareness in solar powered transportation Hans Tholstrup decided to set up the first edition of the World Solar Challenge in 1987. It was a 3000 km race across the Australian outback where competitors from industry research groups and top universities around the globe were invited to compete. General Motors ended up winning the event by a significant margin with their Sunraycer vehicle that achieved speeds of over 40 mph. Contrary to popular belief however solar powered cars are one of the oldest alternative energy vehicles.
Current solar vehicles harness energy from the Sun via Solar panels which are a collected group of solar cells working in tandem towards a common goal. These solid-state devices use quantum mechanical transitions in order to convert a given amount of solar power into electrical power. The electricity produced as a result is then stored in the vehicle's battery in order to run the motor of the vehicle. Batteries in solar-powered vehicles differ from those in standard ICE cars because they are fashioned in a way to impart more power towards the electrical components of the vehicle for a longer duration.
Cells, modules, panels and systems
Multiple solar cells in an integrated group, all oriented in one plane, constitute a solar photovoltaic panel or module. Photovoltaic modules often have a sheet of glass on the sun-facing side, allowing light to pass while protecting the semiconductor wafers. Solar cells are usually connected in series creating additive voltage. Connecting cells in parallel yields a higher current.
However, problems in paralleled cells such as shadow effects can shut down the weaker (less illuminated) parallel string (a number of series connected cells) causing substantial power loss and possible damage because of the reverse bias applied to the shadowed cells by their illuminated partners.
Although modules can be interconnected to create an array with the desired peak DC voltage and loading current capacity, which can be done with or without using independent MPPTs (maximum power point trackers) or, specific to each module, with or without module level power electronic (MLPE) units such as microinverters or DC-DC optimizers. Shunt diodes can reduce shadowing power loss in arrays with series/parallel connected cells.
By 2020, the United States cost per watt for a utility scale system had declined to $0.94.
History
The photovoltaic effect was experimentally demonstrated first by French physicist Edmond Becquerel. In 1839, at age 19, he built the world's first photovoltaic cell in his father's laboratory. Willoughby Smith first described the "Effect of Light on Selenium during the passage of an Electric Current" in a 20 February 1873 issue of Nature. In 1883 Charles Fritts built the first solid state photovoltaic cell by coating the semiconductor selenium with a thin layer of gold to form the junctions; the device was only around 1% efficient. Other milestones include:
1888 – Russian physicist Aleksandr Stoletov built the first cell based on the outer photoelectric effect discovered by Heinrich Hertz in 1887.
1904 – Julius Elster, together with Hans Friedrich Geitel, devised the first practical photoelectric cell.
1905 – Albert Einstein proposed a new quantum theory of light and explained the photoelectric effect in a landmark paper, for which he received the Nobel Prize in Physics in 1921.
1941 – Vadim Lashkaryov discovered p–n junctions in Cu2O and Ag2S protocells.
1946 – Russell Ohl patented the modern junction semiconductor solar cell, while working on the series of advances that would lead to the transistor.
1948 - Introduction to the World of Semiconductors states Kurt Lehovec may have been the first to explain the photo-voltaic effect in the peer reviewed journal Physical Review.
1954 – The first practical photovoltaic cell was publicly demonstrated at Bell Laboratories. The inventors were Calvin Souther Fuller, Daryl Chapin and Gerald Pearson.
1958 – Solar cells gained prominence with their incorporation onto the Vanguard I satellite.
Space applications
Solar cells were first used in a prominent application when they were proposed and flown on the Vanguard satellite in 1958, as an alternative power source to the primary battery power source. By adding cells to the outside of the body, the mission time could be extended with no major changes to the spacecraft or its power systems. In 1959 the United States launched Explorer 6, featuring large wing-shaped solar arrays, which became a common feature in satellites. These arrays consisted of 9600 Hoffman solar cells.
By the 1960s, solar cells were (and still are) the main power source for most Earth orbiting satellites and a number of probes into the solar system, since they offered the best power-to-weight ratio. However, this success was possible because in the space application, power system costs could be high, because space users had few other power options, and were willing to pay for the best possible cells. The space power market drove the development of higher efficiencies in solar cells up until the National Science Foundation "Research Applied to National Needs" program began to push development of solar cells for terrestrial applications.
In the early 1990s the technology used for space solar cells diverged from the silicon technology used for terrestrial panels, with the spacecraft application shifting to gallium arsenide-based III-V semiconductor materials, which then evolved into the modern III-V multijunction photovoltaic cell used on spacecraft.
In recent years, research has moved towards designing and manufacturing lightweight, flexible, and highly efficient solar cells. Terrestrial solar cell technology generally uses photovoltaic cells that are laminated with a layer of glass for strength and protection. Space applications for solar cells require that the cells and arrays are both highly efficient and extremely lightweight. Some newer technology implemented on satellites are multi-junction photovoltaic cells, which are composed of different p–n junctions with varying bandgaps in order to utilize a wider spectrum of the sun's energy. Additionally, large satellites require the use of large solar arrays to produce electricity. These solar arrays need to be broken down to fit in the geometric constraints of the launch vehicle the satellite travels on before being injected into orbit. Historically, solar cells on satellites consisted of several small terrestrial panels folded together. These small panels would be unfolded into a large panel after the satellite is deployed in its orbit. Newer satellites aim to use flexible rollable solar arrays that are very lightweight and can be packed into a very small volume. The smaller size and weight of these flexible arrays drastically decreases the overall cost of launching a satellite due to the direct relationship between payload weight and launch cost of a launch vehicle.
In 2020, the US Naval Research Laboratory conducted its first test of solar power generation in a satellite, the Photovoltaic Radio-frequency Antenna Module (PRAM) experiment aboard the Boeing X-37.
Improved manufacturing methods
Improvements were gradual over the 1960s. This was also the reason that costs remained high, because space users were willing to pay for the best possible cells, leaving no reason to invest in lower-cost, less-efficient solutions. The price was determined largely by the semiconductor industry; their move to integrated circuits in the 1960s led to the availability of larger boules at lower relative prices. As their price fell, the price of the resulting cells did as well. These effects lowered 1971 cell costs to some $100 per watt.
In late 1969 Elliot Berman joined Exxon's task force which was looking for projects 30 years in the future and in April 1973 he founded Solar Power Corporation (SPC), a wholly owned subsidiary of Exxon at that time. The group had concluded that electrical power would be much more expensive by 2000, and felt that this increase in price would make alternative energy sources more attractive. He conducted a market study and concluded that a price per watt of about $20/watt would create significant demand. The team eliminated the steps of polishing the wafers and coating them with an anti-reflective layer, relying on the rough-sawn wafer surface. The team also replaced the expensive materials and hand wiring used in space applications with a printed circuit board on the back, acrylic plastic on the front, and silicone glue between the two, "potting" the cells. Solar cells could be made using cast-off material from the electronics market. By 1973 they announced a product, and SPC convinced Tideland Signal to use its panels to power navigational buoys, initially for the U.S. Coast Guard.
Research and industrial production
Research into solar power for terrestrial applications became prominent with the U.S. National Science Foundation's Advanced Solar Energy Research and Development Division within the "Research Applied to National Needs" program, which ran from 1969 to 1977, and funded research on developing solar power for ground electrical power systems. A 1973 conference, the "Cherry Hill Conference", set forth the technology goals required to achieve this goal and outlined an ambitious project for achieving them, kicking off an applied research program that would be ongoing for several decades. The program was eventually taken over by the Energy Research and Development Administration (ERDA), which was later merged into the U.S. Department of Energy.
Following the 1973 oil crisis, oil companies used their higher profits to start (or buy) solar firms, and were for decades the largest producers. Exxon, ARCO, Shell, Amoco (later purchased by BP) and Mobil all had major solar divisions during the 1970s and 1980s. Technology companies also participated, including General Electric, Motorola, IBM, Tyco and RCA.
Declining costs and exponential growth
Adjusting for inflation, it cost $96 per watt for a solar module in the mid-1970s. Process improvements and a very large boost in production have brought that figure down more than 99%, to 30¢ per watt in 2018
and as low as 20¢ per watt in 2020.
Swanson's law is an observation similar to Moore's Law that states that solar cell prices fall 20% for every doubling of industry capacity. It was featured in an article in the British weekly newspaper The Economist in late 2012. Balance of system costs were then higher than those of the panels. Large commercial arrays could be built, as of 2018, at below $1.00 a watt, fully commissioned.
As the semiconductor industry moved to ever-larger boules, older equipment became inexpensive. Cell sizes grew as equipment became available on the surplus market; ARCO Solar's original panels used cells in diameter. Panels in the 1990s and early 2000s generally used 125 mm wafers; since 2008, almost all new panels use greater than 156mm cells, and by 2020 even larger 182mm ‘M10’ cells. The widespread introduction of flat screen televisions in the late 1990s and early 2000s led to the wide availability of large, high-quality glass sheets to cover the panels.
During the 1990s, polysilicon ("poly") cells became increasingly popular. These cells offer less efficiency than their monosilicon ("mono") counterparts, but they are grown in large vats that reduce cost. By the mid-2000s, poly was dominant in the low-cost panel market, but more recently the mono returned to widespread use.
Manufacturers of wafer-based cells responded to high silicon prices in 2004–2008 with rapid reductions in silicon consumption. In 2008, according to Jef Poortmans, director of IMEC's organic and solar department, current cells use of silicon per watt of power generation, with wafer thicknesses in the neighborhood of 200 microns. Crystalline silicon panels dominate worldwide markets and are mostly manufactured in China and Taiwan. By late 2011, a drop in European demand dropped prices for crystalline solar modules to about $1.09 per watt down sharply from 2010. Prices continued to fall in 2012, reaching $0.62/watt by 4Q2012.
Solar PV is growing fastest in Asia, with China and Japan currently accounting for half of worldwide deployment. Global installed PV capacity reached at least 301 gigawatts in 2016, and grew to supply 1.3% of global power by 2016.
It was anticipated that electricity from PV will be competitive with wholesale electricity costs all across Europe and the energy payback time of crystalline silicon modules can be reduced to below 0.5 years by 2020.
Falling costs are considered one of the biggest factors in the rapid growth of renewable energy, with the cost of solar photovoltaic electricity falling by ~85% between 2010 (when solar and wind made up 1.7% of global electricity generation) and 2021 (where they made up 8.7%). In 2019 solar cells accounted for ~3 % of the world's electricity generation.
Subsidies and grid parity
Solar-specific feed-in tariffs vary by country and within countries. Such tariffs encourage the development of solar power projects. Widespread grid parity, the point at which photovoltaic electricity is equal to or cheaper than grid power without subsidies, likely requires advances on all three fronts. Proponents of solar hope to achieve grid parity first in areas with abundant sun and high electricity costs such as in California and Japan. In 2007 BP claimed grid parity for Hawaii and other islands that otherwise use diesel fuel to produce electricity. George W. Bush set 2015 as the date for grid parity in the US. The Photovoltaic Association reported in 2012 that Australia had reached grid parity (ignoring feed in tariffs).
The price of solar panels fell steadily for 40 years, interrupted in 2004 when high subsidies in Germany drastically increased demand there and greatly increased the price of purified silicon (which is used in computer chips as well as solar panels). The recession of 2008 and the onset of Chinese manufacturing caused prices to resume their decline. In the four years after January 2008 prices for solar modules in Germany dropped from €3 to €1 per peak watt. During that same time production capacity surged with an annual growth of more than 50%. China increased market share from 8% in 2008 to over 55% in the last quarter of 2010. In December 2012 the price of Chinese solar panels had dropped to $0.60/Wp (crystalline modules). (The abbreviation Wp stands for watt peak capacity, or the maximum capacity under optimal conditions.)
As of the end of 2016, it was reported that spot prices for assembled solar panels (not cells) had fallen to a record-low of US$0.36/Wp. The second largest supplier, Canadian Solar Inc., had reported costs of US$0.37/Wp in the third quarter of 2016, having dropped $0.02 from the previous quarter, and hence was probably still at least breaking even. Many producers expected costs would drop to the vicinity of $0.30 by the end of 2017. It was also reported that new solar installations were cheaper than coal-based thermal power plants in some regions of the world, and this was expected to be the case in most of the world within a decade.
Theory
A solar cell is made of semiconducting materials, such as silicon, that have been fabricated into a p–n junction. Such junctions are made by doping one side of the device p-type and the other n-type, for example in the case of silicon by introducing small concentrations of boron or phosphorus respectively.
In operation, photons in sunlight hit the solar cell and are absorbed by the semiconductor. When the photons are absorbed, electrons are excited from the valence band to the conduction band (or from occupied to unoccupied molecular orbitals in the case of an organic solar cell), producing electron-hole pairs.
If the electron-hole pairs are created near the junction between p-type and n-type materials the local electric field sweeps them apart to opposite electrodes, producing an excess of electrons on one side and an excess of holes on the other. When the solar cell is unconnected (or the external electrical load is very high) the electrons and holes will ultimately restore equilibrium by diffusing back across the junction against the field and recombine with each other giving off heat, but if the load is small enough then it is easier for equilibrium to be restored by the excess electrons going around the external circuit, doing useful work along the way.
An array of solar cells converts solar energy into a usable amount of direct current (DC) electricity. An inverter can convert the power to alternating current (AC).
The most commonly known solar cell is configured as a large-area p–n junction made from silicon. Other possible solar cell types are organic solar cells, dye sensitized solar cells, perovskite solar cells, quantum dot solar cells etc. The illuminated side of a solar cell generally has a transparent conducting film for allowing light to enter into the active material and to collect the generated charge carriers. Typically, films with high transmittance and high electrical conductance such as indium tin oxide, conducting polymers or conducting nanowire networks are used for the purpose.
Efficiency
Solar cell efficiency may be broken down into reflectance efficiency, thermodynamic efficiency, charge carrier separation efficiency and conductive efficiency. The overall efficiency is the product of these individual metrics.
The power conversion efficiency of a solar cell is a parameter which is defined by the fraction of incident power converted into electricity.
A solar cell has a voltage dependent efficiency curve, temperature coefficients, and allowable shadow angles.
Due to the difficulty in measuring these parameters directly, other parameters are substituted: thermodynamic efficiency, quantum efficiency, integrated quantum efficiency, VOC ratio, and fill factor. Reflectance losses are a portion of quantum efficiency under "external quantum efficiency". Recombination losses make up another portion of quantum efficiency, VOC ratio, and fill factor. Resistive losses are predominantly categorized under fill factor, but also make up minor portions of quantum efficiency, VOC ratio.
The fill factor is the ratio of the actual maximum obtainable power to the product of the open-circuit voltage and short-circuit current. This is a key parameter in evaluating performance. In 2009, typical commercial solar cells had a fill factor > 0.70. Grade B cells were usually between 0.4 and 0.7. Cells with a high fill factor have a low equivalent series resistance and a high equivalent shunt resistance, so less of the current produced by the cell is dissipated in internal losses.
Single p–n junction crystalline silicon devices are now approaching the theoretical limiting power efficiency of 33.16%, noted as the Shockley–Queisser limit in 1961. In the extreme, with an infinite number of layers, the corresponding limit is 86% using concentrated sunlight.
In 2014, three companies broke the record of 25.6% for a silicon solar cell. Panasonic's was the most efficient. The company moved the front contacts to the rear of the panel, eliminating shaded areas. In addition they applied thin silicon films to the (high quality silicon) wafer's front and back to eliminate defects at or near the wafer surface.
In 2015, a 4-junction GaInP/GaAs//GaInAsP/GaInAs solar cell achieved a new laboratory record efficiency of 46.1% (concentration ratio of sunlight = 312) in a French-German collaboration between the Fraunhofer Institute for Solar Energy Systems (Fraunhofer ISE), CEA-LETI and SOITEC.
In September 2015, Fraunhofer ISE announced the achievement of an efficiency above 20% for epitaxial wafer cells. The work on optimizing the atmospheric-pressure chemical vapor deposition (APCVD) in-line production chain was done in collaboration with NexWafe GmbH, a company spun off from Fraunhofer ISE to commercialize production.
For triple-junction thin-film solar cells, the world record is 13.6%, set in June 2015.
In 2016, researchers at Fraunhofer ISE announced a GaInP/GaAs/Si triple-junction solar cell with two terminals reaching 30.2% efficiency without concentration.
In 2017, a team of researchers at National Renewable Energy Laboratory (NREL), EPFL and CSEM (Switzerland) reported record one-sun efficiencies of 32.8% for dual-junction GaInP/GaAs solar cell devices. In addition, the dual-junction device was mechanically stacked with a Si solar cell, to achieve a record one-sun efficiency of 35.9% for triple-junction solar cells.
Materials
Solar cells are typically named after the semiconducting material they are made of. These materials must have certain characteristics in order to absorb sunlight. Some cells are designed to handle sunlight that reaches the Earth's surface, while others are optimized for use in space. Solar cells can be made of a single layer of light-absorbing material (single-junction) or use multiple physical configurations (multi-junctions) to take advantage of various absorption and charge separation mechanisms.
Solar cells can be classified into first, second and third generation cells. The first generation cells—also called conventional, traditional or wafer-based cells—are made of crystalline silicon, the commercially predominant PV technology, that includes materials such as polysilicon and monocrystalline silicon. Second generation cells are thin film solar cells, that include amorphous silicon, CdTe and CIGS cells and are commercially significant in utility-scale photovoltaic power stations, building integrated photovoltaics or in small stand-alone power system. The third generation of solar cells includes a number of thin-film technologies often described as emerging photovoltaics—most of them have not yet been commercially applied and are still in the research or development phase. Many use organic materials, often organometallic compounds as well as inorganic substances. Despite the fact that their efficiencies had been low and the stability of the absorber material was often too short for commercial applications, there is research into these technologies as they promise to achieve the goal of producing low-cost, high-efficiency solar cells. As of 2016, the most popular and efficient solar cells were those made from thin wafers of silicon which are also the oldest solar cell technology.
Crystalline silicon
By far, the most prevalent bulk material for solar cells is crystalline silicon (c-Si), also known as "solar grade silicon". Bulk silicon is separated into multiple categories according to crystallinity and crystal size in the resulting ingot, ribbon or wafer. These cells are entirely based around the concept of a p–n junction. Solar cells made of c-Si are made from wafers between 160 and 240 micrometers thick.
Monocrystalline silicon
Monocrystalline silicon (mono-Si) solar cells feature a single-crystal composition that enables electrons to move more freely than in a multi-crystal configuration. Consequently, monocrystalline solar panels deliver a higher efficiency than their multicrystalline counterparts. The corners of the cells look clipped, like an octagon, because the wafer material is cut from cylindrical ingots, that are typically grown by the Czochralski process. Solar panels using mono-Si cells display a distinctive pattern of small white diamonds.
Epitaxial silicon development
Epitaxial wafers of crystalline silicon can be grown on a monocrystalline silicon "seed" wafer by chemical vapor deposition (CVD), and then detached as self-supporting wafers of some standard thickness (e.g., 250 μm) that can be manipulated by hand, and directly substituted for wafer cells cut from monocrystalline silicon ingots. Solar cells made with this "kerfless" technique can have efficiencies approaching those of wafer-cut cells, but at appreciably lower cost if the CVD can be done at atmospheric pressure in a high-throughput inline process. The surface of epitaxial wafers may be textured to enhance light absorption.
In June 2015, it was reported that heterojunction solar cells grown epitaxially on n-type monocrystalline silicon wafers had reached an efficiency of 22.5% over a total cell area of 243.4 cm.
Polycrystalline silicon
Polycrystalline silicon, or multicrystalline silicon (multi-Si) cells are made from cast square ingots—large blocks of molten silicon carefully cooled and solidified. They consist of small crystals giving the material its typical metal flake effect. Polysilicon cells are the most common type used in photovoltaics and are less expensive, but also less efficient, than those made from monocrystalline silicon.
Ribbon silicon
Ribbon silicon is a type of polycrystalline silicon—it is formed by drawing flat thin films from molten silicon and results in a polycrystalline structure. These cells are cheaper to make than multi-Si, due to a great reduction in silicon waste, as this approach does not require sawing from ingots. However, they are also less efficient.
Mono-like-multi silicon (MLM)
This form was developed in the 2000s and introduced commercially around 2009. Also called cast-mono, this design uses polycrystalline casting chambers with small "seeds" of mono material. The result is a bulk mono-like material that is polycrystalline around the outsides. When sliced for processing, the inner sections are high-efficiency mono-like cells (but square instead of "clipped"), while the outer edges are sold as conventional poly. This production method results in mono-like cells at poly-like prices.
Thin film
Thin-film technologies reduce the amount of active material in a cell. Most designs sandwich active material between two panes of glass. Since silicon solar panels only use one pane of glass, thin film panels are approximately twice as heavy as crystalline silicon panels, although they have a smaller ecological impact (determined from life cycle analysis).
Cadmium telluride
Cadmium telluride is the only thin film material so far to rival crystalline silicon in cost/watt. However cadmium is highly toxic and tellurium (anion: "telluride") supplies are limited. The cadmium present in the cells would be toxic if released. However, release is impossible during normal operation of the cells and is unlikely during fires in residential roofs. A square meter of CdTe contains approximately the same amount of Cd as a single C cell nickel-cadmium battery, in a more stable and less soluble form.
Copper indium gallium selenide
Copper indium gallium selenide (CIGS) is a direct band gap material. It has the highest efficiency (~20%) among all commercially significant thin film materials (see CIGS solar cell). Traditional methods of fabrication involve vacuum processes including co-evaporation and sputtering. Recent developments at IBM and Nanosolar attempt to lower the cost by using non-vacuum solution processes.
Silicon thin film
Silicon thin-film cells are mainly deposited by chemical vapor deposition (typically plasma-enhanced, PE-CVD) from silane gas and hydrogen gas. Depending on the deposition parameters, this can yield amorphous silicon (a-Si or a-Si:H), protocrystalline silicon or nanocrystalline silicon (nc-Si or nc-Si:H), also called microcrystalline silicon.
Amorphous silicon is the most well-developed thin film technology to-date. An amorphous silicon (a-Si) solar cell is made of non-crystalline or microcrystalline silicon. Amorphous silicon has a higher bandgap (1.7 eV) than crystalline silicon (c-Si) (1.1 eV), which means it absorbs the visible part of the solar spectrum more strongly than the higher power density infrared portion of the spectrum. The production of a-Si thin film solar cells uses glass as a substrate and deposits a very thin layer of silicon by plasma-enhanced chemical vapor deposition (PECVD).
Protocrystalline silicon with a low volume fraction of nanocrystalline silicon is optimal for high open-circuit voltage. Nc-Si has about the same bandgap as c-Si and nc-Si and a-Si can advantageously be combined in thin layers, creating a layered cell called a tandem cell. The top cell in a-Si absorbs the visible light and leaves the infrared part of the spectrum for the bottom cell in nc-Si.
Gallium arsenide thin film
The semiconductor material gallium arsenide (GaAs) is also used for single-crystalline thin film solar cells. Although GaAs cells are very expensive, they hold the world's record in efficiency for a single-junction solar cell at 28.8%. Typically fabricated on crystalline silicon wafer with a 41% fill factor, by moving to porous silicon fill factor can be increased to 56% with potentially reduced cost. Using less active GaAs material by fabricating nanowires is another potential pathway to cost reduction. GaAs is more commonly used in multijunction photovoltaic cells for concentrated photovoltaics (CPV, HCPV) and for solar panels on spacecraft, as the industry favours efficiency over cost for space-based solar power. Based on the previous literature and some theoretical analysis, there are several reasons why GaAs has such high power conversion efficiency. First, GaAs bandgap is 1.43ev which is almost ideal for solar cells. Second, because Gallium is a by-product of the smelting of other metals, GaAs cells are relatively insensitive to heat and it can keep high efficiency when temperature is quite high. Third, GaAs has the wide range of design options. Using GaAs as active layer in solar cell, engineers can have multiple choices of other layers which can better generate electrons and holes in GaAs.
Multijunction cells
Multi-junction cells consist of multiple thin films, each essentially a solar cell grown on top of another, typically using metalorganic vapour phase epitaxy. Each layer has a different band gap energy to allow it to absorb electromagnetic radiation over a different portion of the spectrum. Multi-junction cells were originally developed for special applications such as satellites and space exploration, but are now used increasingly in terrestrial concentrator photovoltaics (CPV), an emerging technology that uses lenses and curved mirrors to concentrate sunlight onto small, highly efficient multi-junction solar cells. By concentrating sunlight up to a thousand times, High concentration photovoltaics (HCPV) has the potential to outcompete conventional solar PV in the future.
Tandem solar cells based on monolithic, series connected, gallium indium phosphide (GaInP), gallium arsenide (GaAs), and germanium (Ge) p–n junctions, are increasing sales, despite cost pressures. Between December 2006 and December 2007, the cost of 4N gallium metal rose from about $350 per kg to $680 per kg. Additionally, germanium metal prices have risen substantially to $1000–1200 per kg this year. Those materials include gallium (4N, 6N and 7N Ga), arsenic (4N, 6N and 7N) and germanium, pyrolitic boron nitride (pBN) crucibles for growing crystals, and boron oxide, these products are critical to the entire substrate manufacturing industry.
A triple-junction cell, for example, may consist of the semiconductors: GaAs, Ge, and . Triple-junction GaAs solar cells were used as the power source of the Dutch four-time World Solar Challenge winners Nuna in 2003, 2005 and 2007 and by the Dutch solar cars Solutra (2005), Twente One (2007) and 21Revolution (2009). GaAs based multi-junction devices are the most efficient solar cells to date. On 15 October 2012, triple junction metamorphic cells reached a record high of 44%. In 2022, researchers at Fraunhofer Institute for Solar Energy Systems ISE in Freiburg, Germany, demonstrated a record solar cell efficiency of 47.6% under 665-fold sunlight concentration with a four-junction concentrator solar cell.
GaInP/Si dual-junction solar cells
In 2016, a new approach was described for producing hybrid photovoltaic wafers combining the high efficiency of III-V multi-junction solar cells with the economies and wealth of experience associated with silicon. The technical complications involved in growing the III-V material on silicon at the required high temperatures, a subject of study for some 30 years, are avoided by epitaxial growth of silicon on GaAs at low temperature by plasma-enhanced chemical vapor deposition (PECVD).
Si single-junction solar cells have been widely studied for decades and are reaching their practical efficiency of ~26% under 1-sun conditions. Increasing this efficiency may require adding more cells with bandgap energy larger than 1.1 eV to the Si cell, allowing to convert short-wavelength photons for generation of additional voltage. A dual-junction solar cell with a band gap of 1.6–1.8 eV as a top cell can reduce thermalization loss, produce a high external radiative efficiency and achieve theoretical efficiencies over 45%. A tandem cell can be fabricated by growing the GaInP and Si cells. Growing them separately can overcome the 4% lattice constant mismatch between Si and the most common III–V layers that prevent direct integration into one cell. The two cells therefore are separated by a transparent glass slide so the lattice mismatch does not cause strain to the system. This creates a cell with four electrical contacts and two junctions that demonstrated an efficiency of 18.1%. With a fill factor (FF) of 76.2%, the Si bottom cell reaches an efficiency of 11.7% (± 0.4) in the tandem device, resulting in a cumulative tandem cell efficiency of 29.8%. This efficiency exceeds the theoretical limit of 29.4% and the record experimental efficiency value of a Si 1-sun solar cell, and is also higher than the record-efficiency 1-sun GaAs device. However, using a GaAs substrate is expensive and not practical. Hence researchers try to make a cell with two electrical contact points and one junction, which does not need a GaAs substrate. This means there will be direct integration of GaInP and Si.
Research in solar cells
Perovskite solar cells
Perovskite solar cells are solar cells that include a perovskite-structured material as the active layer. Most commonly, this is a solution-processed hybrid organic-inorganic tin or lead halide based material. Efficiencies have increased from below 5% at their first usage in 2009 to 25.5% in 2020, making them a very rapidly advancing technology and a hot topic in the solar cell field. Researchers at University of Rochester reported in 2023 that significant further improvements in cell efficiency can be achieved by utilizing Purcell effect.
Perovskite solar cells are also forecast to be extremely cheap to scale up, making them a very attractive option for commercialisation. So far most types of perovskite solar cells have not reached sufficient operational stability to be commercialised, although many research groups are investigating ways to solve this. Energy and environmental sustainability of perovskite solar cells and tandem perovskite are shown to be dependent on the structures. Photonic front contacts for light management can improve the perovskite cells' performance, via enhanced broadband absorption, while allowing better operational stability due to protection against the harmful high-energy (above Visible) radiation. The inclusion of the toxic element lead in the most efficient perovskite solar cells is a potential problem for commercialisation.
Bifacial solar cells
With a transparent rear side, bifacial solar cells can absorb light from both the front and rear sides. Hence, they can produce more electricity than conventional monofacial solar cells. The first patent of bifacial solar cells was filed by Japanese researcher Hiroshi Mori, in 1966. Later, it is said that Russia was the first to deploy bifacial solar cells in their space program in the 1970s. In 1976, the Institute for Solar Energy of the Technical University of Madrid, began a research program for the development of bifacial solar cells led by Prof. Antonio Luque. Based on 1977 US and Spanish patents by Luque, a practical bifacial cell was proposed with a front face as anode and a rear face as cathode; in previously reported proposals and attempts both faces were anodic and interconnection between cells was complicated and expensive. In 1980, Andrés Cuevas, a PhD student in Luque's team, demonstrated experimentally a 50% increase in output power of bifacial solar cells, relative to identically oriented and tilted monofacial ones, when a white background was provided. In 1981 the company Isofoton was founded in Málaga to produce the developed bifacial cells, thus becoming the first industrialization of this PV cell technology. With an initial production capacity of 300 kW/yr of bifacial solar cells, early landmarks of Isofoton's production were the 20kWp power plant in San Agustín de Guadalix, built in 1986 for Iberdrola, and an off grid installation by 1988 also of 20kWp in the village of Noto Gouye Diama (Senegal) funded by the Spanish international aid and cooperation programs.
Due to the reduced manufacturing cost, companies have again started to produce commercial bifacial modules since 2010. By 2017, there were at least eight certified PV manufacturers providing bifacial modules in North America. The International Technology Roadmap for Photovoltaics (ITRPV) predicted that the global market share of bifacial technology will expand from less than 5% in 2016 to 30% in 2027.
Due to the significant interest in the bifacial technology, a recent study has investigated the performance and optimization of bifacial solar modules worldwide. The results indicate that, across the globe, ground-mounted bifacial modules can only offer ~10% gain in annual electricity yields compared to the monofacial counterparts for a ground albedo coefficient of 25% (typical for concrete and vegetation groundcovers). However, the gain can be increased to ~30% by elevating the module 1 m above the ground and enhancing the ground albedo coefficient to 50%. Sun et al. also derived a set of empirical equations that can optimize bifacial solar modules analytically. In addition, there is evidence that bifacial panels work better than traditional panels in snowy environments as bifacials on dual-axis trackers made 14% more electricity in a year than their monofacial counterparts and 40% during the peak winter months.
An online simulation tool is available to model the performance of bifacial modules in any arbitrary location across the entire world. It can also optimize bifacial modules as a function of tilt angle, azimuth angle, and elevation above the ground.
Intermediate band
Intermediate band photovoltaics in solar cell research provides methods for exceeding the Shockley–Queisser limit on the efficiency of a cell. It introduces an intermediate band (IB) energy level in between the valence and conduction bands. Theoretically, introducing an IB allows two photons with energy less than the bandgap to excite an electron from the valence band to the conduction band. This increases the induced photocurrent and thereby efficiency.
Luque and Marti first derived a theoretical limit for an IB device with one midgap energy level using detailed balance. They assumed no carriers were collected at the IB and that the device was under full concentration. They found the maximum efficiency to be 63.2%, for a bandgap of 1.95eV with the IB 0.71eV from either the valence or conduction band.
Under one sun illumination the limiting efficiency is 47%. Several means are under study to realize IB semiconductors with such optimum 3-bandgap configuration, namely via materials engineering (controlled inclusion of deep level impurities or highly-mismatched alloys) and nano-structuring (quantum-dots in host hetero-crystals).
Liquid inks
In 2014, researchers at California NanoSystems Institute discovered using kesterite and perovskite improved electric power conversion efficiency for solar cells.
In December 2022, it was reported that MIT researchers had developed ultralight fabric solar cells. These cells offer a weight one-hundredth that of traditional panels while generating 18 times more power per kilogram. Thinner than a human hair, these cells can be laminated onto various surfaces, such as boat sails, tents, tarps, or drone wings, to extend their functionality. Using ink-based materials and scalable techniques, researchers coat the solar cell structure with printable electronic inks, completing the module with screen-printed electrodes. Tested on high-strength fabric, the cells produce 370 watts-per-kilogram, representing an improvement over conventional solar cells.
Upconversion and downconversion
Photon upconversion is the process of using two low-energy (e.g., infrared) photons to produce one higher energy photon; downconversion is the process of using one high energy photon (e.g., ultraviolet) to produce two lower energy photons. Either of these techniques could be used to produce higher efficiency solar cells by allowing solar photons to be more efficiently used. The difficulty, however, is that the conversion efficiency of existing phosphors exhibiting up- or down-conversion is low, and is typically narrow band.
One upconversion technique is to incorporate lanthanide-doped materials (, , or a combination), taking advantage of their luminescence to convert infrared radiation to visible light. Upconversion process occurs when two infrared photons are absorbed by rare-earth ions to generate a (high-energy) absorbable photon. As example, the energy transfer upconversion process (ETU), consists in successive transfer processes between excited ions in the near infrared. The upconverter material could be placed below the solar cell to absorb the infrared light that passes through the silicon. Useful ions are most commonly found in the trivalent state. ions have been the most used. ions absorb solar radiation around 1.54 μm. Two ions that have absorbed this radiation can interact with each other through an upconversion process. The excited ion emits light above the Si bandgap that is absorbed by the solar cell and creates an additional electron–hole pair that can generate current. However, the increased efficiency was small. In addition, fluoroindate glasses have low phonon energy and have been proposed as suitable matrix doped with ions.
Light-absorbing dyes
Dye-sensitized solar cells (DSSCs) are made of low-cost materials and do not need elaborate manufacturing equipment, so they can be made in a DIY fashion. In bulk it should be significantly less expensive than older solid-state cell designs. DSSC's can be engineered into flexible sheets and although its conversion efficiency is less than the best thin film cells, its price/performance ratio may be high enough to allow them to compete with fossil fuel electrical generation.
Typically a ruthenium metalorganic dye (Ru-centered) is used as a monolayer of light-absorbing material, which is adsorbed onto a thin film of titanium dioxide. The dye-sensitized solar cell depends on this mesoporous layer of nanoparticulate titanium dioxide (TiO2) to greatly amplify the surface area (200–300 m2/g , as compared to approximately 10 m2/g of flat single crystal) which allows for a greater number of dyes per solar cell area (which in term in increases the current). The photogenerated electrons from the light absorbing dye are passed on to the n-type and the holes are absorbed by an electrolyte on the other side of the dye. The circuit is completed by a redox couple in the electrolyte, which can be liquid or solid. This type of cell allows more flexible use of materials and is typically manufactured by screen printing or ultrasonic nozzles, with the potential for lower processing costs than those used for bulk solar cells. However, the dyes in these cells also suffer from degradation under heat and UV light and the cell casing is difficult to seal due to the solvents used in assembly. Due to this reason, researchers have developed solid-state dye-sensitized solar cells that use a solid electrolyte to avoid leakage. The first commercial shipment of DSSC solar modules occurred in July 2009 from G24i Innovations.
Quantum dots
Quantum dot solar cells (QDSCs) are based on the Gratzel cell, or dye-sensitized solar cell architecture, but employ low band gap semiconductor nanoparticles, fabricated with crystallite sizes small enough to form quantum dots (such as CdS, CdSe, , PbS, etc.), instead of organic or organometallic dyes as light absorbers. Due to the toxicity associated with Cd and Pb based compounds there are also a series of "green" QD sensitizing materials in development (such as CuInS2, CuInSe2 and CuInSeS). QD's size quantization allows for the band gap to be tuned by simply changing particle size. They also have high extinction coefficients and have shown the possibility of multiple exciton generation.
In a QDSC, a mesoporous layer of titanium dioxide nanoparticles forms the backbone of the cell, much like in a DSSC. This layer can then be made photoactive by coating with semiconductor quantum dots using chemical bath deposition, electrophoretic deposition or successive ionic layer adsorption and reaction. The electrical circuit is then completed through the use of a liquid or solid redox couple. The efficiency of QDSCs has increased to over 5% shown for both liquid-junction and solid state cells, with a reported peak efficiency of 11.91%. In an effort to decrease production costs, the Prashant Kamat research group demonstrated a solar paint made with and CdSe that can be applied using a one-step method to any conductive surface with efficiencies over 1%. However, the absorption of quantum dots (QDs) in QDSCs is weak at room temperature. The plasmonic nanoparticles can be utilized to address the weak absorption of QDs (e.g., nanostars). Adding an external infrared pumping source to excite intraband and interband transition of QDs is another solution.
Organic/polymer solar cells
Organic solar cells and polymer solar cells are built from thin films (typically 100 nm) of organic semiconductors including polymers, such as polyphenylene vinylene and small-molecule compounds like copper phthalocyanine (a blue or green organic pigment) and carbon fullerenes and fullerene derivatives such as PCBM.
They can be processed from liquid solution, offering the possibility of a simple roll-to-roll printing process, potentially leading to inexpensive, large-scale production. In addition, these cells could be beneficial for some applications where mechanical flexibility and disposability are important. Current cell efficiencies are, however, very low, and practical devices are essentially non-existent.
Energy conversion efficiencies achieved to date using conductive polymers are very low compared to inorganic materials. However, Konarka Power Plastic reached efficiency of 8.3% and organic tandem cells in 2012 reached 11.1%.
The active region of an organic device consists of two materials, one electron donor and one electron acceptor. When a photon is converted into an electron hole pair, typically in the donor material, the charges tend to remain bound in the form of an exciton, separating when the exciton diffuses to the donor-acceptor interface, unlike most other solar cell types. The short exciton diffusion lengths of most polymer systems tend to limit the efficiency of such devices. Nanostructured interfaces, sometimes in the form of bulk heterojunctions, can improve performance.
In 2011, MIT and Michigan State researchers developed solar cells with a power efficiency close to 2% with a transparency to the human eye greater than 65%, achieved by selectively absorbing the ultraviolet and near-infrared parts of the spectrum with small-molecule compounds. Researchers at UCLA more recently developed an analogous polymer solar cell, following the same approach, that is 70% transparent and has a 4% power conversion efficiency. These lightweight, flexible cells can be produced in bulk at a low cost and could be used to create power generating windows.
In 2013, researchers announced polymer cells with some 3% efficiency. They used block copolymers, self-assembling organic materials that arrange themselves into distinct layers. The research focused on P3HT-b-PFTBT that separates into bands some 16 nanometers wide.
Adaptive cells
Adaptive cells change their absorption/reflection characteristics depending on environmental conditions. An adaptive material responds to the intensity and angle of incident light. At the part of the cell where the light is most intense, the cell surface changes from reflective to adaptive, allowing the light to penetrate the cell. The other parts of the cell remain reflective increasing the retention of the absorbed light within the cell.
In 2014, a system was developed that combined an adaptive surface with a glass substrate that redirect the absorbed to a light absorber on the edges of the sheet. The system also includes an array of fixed lenses/mirrors to concentrate light onto the adaptive surface. As the day continues, the concentrated light moves along the surface of the cell. That surface switches from reflective to adaptive when the light is most concentrated and back to reflective after the light moves along.
Surface texturing
For the past years, researchers have been trying to reduce the price of solar cells while maximizing efficiency. Thin-film solar cell is a cost-effective second generation solar cell with much reduced thickness at the expense of light absorption efficiency. Efforts to maximize light absorption efficiency with reduced thickness have been made. Surface texturing is one of techniques used to reduce optical losses to maximize light absorbed. Currently, surface texturing techniques on silicon photovoltaics are drawing much attention. Surface texturing could be done in multiple ways. Etching single crystalline silicon substrate can produce randomly distributed square based pyramids on the surface using anisotropic etchants. Recent studies show that c-Si wafers could be etched down to form nano-scale inverted pyramids. Multicrystalline silicon solar cells, due to poorer crystallographic quality, are less effective than single crystal solar cells, but mc-Si solar cells are still being used widely due to less manufacturing difficulties. It is reported that multicrystalline solar cells can be surface-textured to yield solar energy conversion efficiency comparable to that of monocrystalline silicon cells, through isotropic etching or photolithography techniques. Incident light rays onto a textured surface do not reflect back out to the air as opposed to rays onto a flat surface. Rather some light rays are bounced back onto the other surface again due to the geometry of the surface. This process significantly improves light to electricity conversion efficiency, due to increased light absorption. This texture effect as well as the interaction with other interfaces in the PV module is a challenging optical simulation task. A particularly efficient method for modeling and optimization is the OPTOS formalism. In 2012, researchers at MIT reported that c-Si films textured with nanoscale inverted pyramids could achieve light absorption comparable to 30 times thicker planar c-Si. In combination with anti-reflective coating, surface texturing technique can effectively trap light rays within a thin film silicon solar cell. Consequently, required thickness for solar cells decreases with the increased absorption of light rays.
Encapsulation
Solar cells are commonly encapsulated in a transparent polymeric resin to protect the delicate solar cell regions for coming into contact with moisture, dirt, ice, and other conditions expected either during operation or when used outdoors. The encapsulants are commonly made from polyvinyl acetate or glass. Most encapsulants are uniform in structure and composition, which increases light collection owing to light trapping from total internal reflection of light within the resin. Research has been conducted into structuring the encapsulant to provide further collection of light. Such encapsulants have included roughened glass surfaces, diffractive elements, prism arrays, air prisms, v-grooves, diffuse elements, as well as multi-directional waveguide arrays. Prism arrays show an overall 5% increase in the total solar energy conversion. Arrays of vertically aligned broadband waveguides provide a 10% increase at normal incidence, as well as wide-angle collection enhancement of up to 4%, with optimized structures yielding up to a 20% increase in short circuit current. Active coatings that convert infrared light into visible light have shown a 30% increase. Nanoparticle coatings inducing plasmonic light scattering increase wide-angle conversion efficiency up to 3%. Optical structures have also been created in encapsulation materials to effectively "cloak" the metallic front contacts.
Autonomous maintenance
Novel self-cleaning mechanisms for solar panels are being developed. For instance, in 2019 via wet-chemically etched nanowires and a hydrophobic coating on the surface water droplets could remove 98% of dust particles, which may be especially relevant for applications in the desert.
In March 2022, MIT researchers announced the development of a waterless cleaning system for solar panels and mirrors to address the issue of dust accumulation, which can reduce solar output by up to 30 percent in one month. This system utilizes electrostatic repulsion to detach dust particles from the panel's surface, eliminating the need for water or brushes. An electrical charge imparted to the dust particles by passing a simple electrode over the panel causes them to be repelled by a charge applied to the panel itself. The system can be automated using a basic electric motor and guide rails.
Manufacture
Solar cells share some of the same processing and manufacturing techniques as other semiconductor devices. However, the strict requirements for cleanliness and quality control of semiconductor fabrication are more relaxed for solar cells, lowering costs.
Polycrystalline silicon wafers are made by wire-sawing block-cast silicon ingots into 180 to 350 micrometer wafers. The wafers are usually lightly p-type-doped. A surface diffusion of n-type dopants is performed on the front side of the wafer. This forms a p–n junction a few hundred nanometers below the surface.
Anti-reflection coatings are then typically applied to increase the amount of light coupled into the solar cell. Silicon nitride has gradually replaced titanium dioxide as the preferred material, because of its excellent surface passivation qualities. It prevents carrier recombination at the cell surface. A layer several hundred nanometers thick is applied using plasma-enhanced chemical vapor deposition. Some solar cells have textured front surfaces that, like anti-reflection coatings, increase the amount of light reaching the wafer. Such surfaces were first applied to single-crystal silicon, followed by multicrystalline silicon somewhat later.
A full area metal contact is made on the back surface, and a grid-like metal contact made up of fine "fingers" and larger "bus bars" are screen-printed onto the front surface using a silver paste. This is an evolution of the so-called "wet" process for applying electrodes, first described in a US patent filed in 1981 by Bayer AG. The rear contact is formed by screen-printing a metal paste, typically aluminium. Usually this contact covers the entire rear, though some designs employ a grid pattern. The paste is then fired at several hundred degrees Celsius to form metal electrodes in ohmic contact with the silicon. Some companies use an additional electroplating step to increase efficiency. After the metal contacts are made, the solar cells are interconnected by flat wires or metal ribbons, and assembled into modules or "solar panels". Solar panels have a sheet of tempered glass on the front, and a polymer encapsulation on the back.
Different types of manufacturing and recycling partly determine how effective it is in decreasing emissions and having a positive environmental effect. Such differences and effectiveness could be quantified for production of the most optimal types of products for different purposes in different regions across time.
Manufacturers and certification
National Renewable Energy Laboratory tests and validates solar technologies. Three reliable groups certify solar equipment: UL and IEEE (both U.S. standards) and IEC.
The IEA's 2022 Special Report highlights China's dominance over the solar PV supply chain, with an investment exceeding USD 50 billion and the creation of around 300,000 jobs since 2011. China commands over 80% of all manufacturing stages for solar panels. This control has drastically cut costs but also led to issues like supply-demand imbalances and polysilicon production constraints. Nevertheless, China's strategic policies have reduced solar PV costs by more than 80%, increasing global affordability. In 2021, China's solar PV exports were over USD 30 billion.
Meeting global energy and climate targets necessitates a major expansion in solar PV manufacturing, aiming for over 630 GW by 2030 according to the IEA's "Roadmap to Net Zero Emissions by 2050". China's dominance, controlling nearly 95% of key solar PV components and 40% of the world's polysilicon production in Xinjiang, poses risks of supply shortages and cost surges. Critical mineral demand, like silver, may exceed 30% of 2020's global production by 2030.
In 2021, China's share of solar PV module production reached approximately 70%, an increase from 50% in 2010. Other key producers included Vietnam (5%), Malaysia (4%), Korea (4%), and Thailand (2%), with much of their production capacity developed by Chinese companies aimed at exports, notably to the United States.
China
As of September 2018, sixty percent of the world's solar photovoltaic modules were made in China. As of May 2018, the largest photovoltaic plant in the world is located in the Tengger desert in China. In 2018, China added more photovoltaic installed capacity (in GW) than the next 9 countries combined. In 2021, China's share of solar PV module production reached approximately 70%.
In the first half of 2023, China's production of PV modules exceeded 220 GW, marking an increase of over 62% compared to the same period in 2022. In 2022, China maintained its position as the world's largest PV module producer, holding a dominant market share of 77.8%.
Vietnam
In 2022, Vietnam was the second-largest PV module producer, only behind China, with its production capacity rising to 24.1 GW, marking a significant 47% increase from the 16.4 GW produced in 2021. Vietnam accounts for 6.4% of the world's photovoltaic production.
Malaysia
In 2022, Malaysia was the third-largest PV module producer, with a production capacity of 10.8 GW, accounting for 2.8% of global production. This placed it behind China, which dominated with 77.8%, and Vietnam, which contributed 6.4%.
United States
Solar energy production in the U.S. has doubled from 2013 to 2019. This was driven first by the falling price of quality silicon, and later simply by the globally plunging cost of photovoltaic modules. In 2018, the U.S. added 10.8GW of installed solar photovoltaic energy, an increase of 21%.
Latin America: Latin America has emerged as a promising region for solar energy development in recent years, with over 10 GW of installations in 2020. The solar market in Latin America has been driven by abundant solar resources, falling costs, competitive auctions and growing electricity demand. Some of the leading countries for solar energy in Latin America are Brazil, Mexico, Chile and Argentina. However, the solar market in Latin America also faces some challenges, such as political instability, financing gaps and power transmission bottlenecks.
Middle East and Africa: The Middle East and Africa has also experienced significant growth in solar energy deployment in recent years, with over 8 GW installations in 2020. The solar market in the Middle East and Africa has been driven by the low-cost generation of solar energy, the diversification of energy sources, the fight against climate change and rural electrification are motivated. Some of the notable countries for solar energy in the Middle East and Africa are Saudi Arabia, United Arab Emirates, Egypt, Morocco and South Africa. However, the solar market in the Middle East and Africa also faces several obstacles, including social unrest, regulatory uncertainty and technical barriers.
Materials sourcing
Like many other energy generation technologies, the manufacture of solar cells, especially its rapid expansion, has many environmental and supply-chain implications. Global mining may adapt and potentially expand for sourcing the needed minerals which vary per type of solar cell. Recycling solar panels could be a source for materials that would otherwise need to be mined.
Disposal
Solar cells degrade over time and lose their efficiency. Solar cells in extreme climates, such as desert or polar, are more prone to degradation due to exposure to harsh UV light and snow loads respectively. Usually, solar panels are given a lifespan of 25–30 years before they get decommissioned.
The International Renewable Energy Agency estimated that the amount of solar panel electronic waste generated in 2016 was 43,500–250,000 metric tons. This number is estimated to increase substantially by 2030, reaching an estimated waste volume of 60–78 million metric tons in 2050.
Recycling
The most widely used solar cells in the market are crystalline solar cells. A product is truly recyclable if it can harvested again. In the 2016 Paris Agreement, 195 countries agreed to reduce their carbon emissions by shifting their focus away from fossil fuels and towards renewable energy sources. Owing to this, Solar will be a major contributor to electricity generation all over the world. So, there will be a plethora of solar panels to be recycled after the end of their life cycle. In fact, many researchers around the globe have voiced their concern about finding ways to use silicon cells after recycling.
Additionally, these cells have hazardous elements/compounds, including lead (Pb), cadmium (Cd) or cadmium sulfide (CdS), selenium (Se), and barium (Ba) as dopants aside from the valuables silicon (Si), aluminum (Al), silver (Ag), and copper (Cu). The harmful elements/compounds if not disposed of with the proper technique can have severe harmful effects on human life and wildlife alike.
There are various ways c-Si can be recycled. Mainly thermal and chemical separation methods are used. This happens in two stages
PV solar cell separation: in thermal delamination, the ethylene vinyl acetate (EVA) is removed and materials such as glass, Tedlar®, aluminium frame, steel, copper and plastics are separated;
cleansing the surface of PV solar cells: unwanted layers (antireflection layer, metal coating and p–n semiconductor) are removed from the silicon solar cells separated from the PV modules; as a result, the silicon substrate, suitable for re-use, can be recovered.
The First Solar panel recycling plant opened in Rousset, France in 2018. It was set to recycle 1300 tonnes of solar panel waste a year, and can increase its capacity to 4000 tonnes. If recycling is driven only by market-based prices, rather than also environmental regulations, the economic incentives for recycling remain uncertain and as of 2021 the environmental impact of different types of developed recycling techniques still need to be quantified.
| Technology | Energy and fuel | null |
2353769 | https://en.wikipedia.org/wiki/Thorianite | Thorianite | Thorianite is a rare thorium oxide mineral, ThO2. It was originally described by Ananda Coomaraswamy in 1904 as uraninite, but recognized as a new species by Wyndham R. Dunstan. It was so named by Dunstan on account of its high percentage of thorium; it also contains the oxides of uranium, lanthanum, cerium, praseodymium and neodymium. Helium is present, and the mineral is slightly less radioactive than pitchblende, but is harder to shield due to its high energy gamma rays. It is common in the alluvial gem-gravels of Sri Lanka, where it occurs mostly as water worn, small, heavy, black, cubic crystals. The largest crystals are usually near 1.5 cm. Larger crystals, up to , have been reported from Madagascar.
Chemistry
Based on color, specific gravity and composition three types of thorianite are distinguished:
α-thorianite
β-thorianite
γ-thorianite
Thorianite and uraninite form a complete solid solution series in synthetic and natural material. The division between the two species is at Th:U = 1:1 with U possibly making up to 46.50% and Th ranging up to 87.88%. Rare earths, chiefly Ce, substitute for Th in amounts up to 8% by weight. Ce is probably present as Ce4+. Complete series is known in synthetic material between CeO2 - PrO2 - ThO2 - UO2. Small amounts of Fe3+ and Zr also may be isomorphous with Th. Pb present is probably radiogenic.
Varieties
Aldanite – a variety of thorianite containing 14.9% to 29.0% UO2 and 11.2% to 12.5% PbO.
Uranothorianite
Thorianite Cerian
Thorianite La bearing
Occurrence
Usually found in alluvial deposits, beach sands, heavy mineral placers, and pegmatites.
Sri Lanka – In stream gravels, Galle district, Southern Province; Balangoda district; near Kodrugala, Sabaragamuwa Province; and from a pegmatite in Bambarabotuwa area.
India – Reported from beach sands of Travancore (Kerala).
Madagascar – Found in alluvial deposits of Betroka and Andolobe. Also as very large crystals from Fort Dauphin; at Andranondambo and other localities.
Russia – In black sands of a gold placer on Boshogoch River, Transbaikalia, Siberia; in the Kovdor Massif by Kovdor, Kola Peninsula; in the Yenisei Range, Siberia.
United States – reported from Easton, Pennsylvania; black sands in Missouri River, near Helena, Montana; Scott River, Siskiyou County, California; black sands in Nixon Fork and Wiseman districts, Alaska.
Canada – Reported with uraninite in a pegmatite on Charlebois Lake, east of Lake Athabasca; Uranon variety reported from pegmatite and metesomatized zones in crystalline limestones from many locations in Quebec and Ontario.
South Africa – Occurs with baddeleyite as an accessory in carbonatite at Phalaborwa, Eastern Transvaal.
Democratic Republic of Congo - Kasaï region
| Physical sciences | Minerals | Earth science |
2354317 | https://en.wikipedia.org/wiki/Sulfolobus | Sulfolobus | Sulfolobus is a genus of microorganism in the family Sulfolobaceae. It belongs to the archaea domain.
Sulfolobus species grow in volcanic springs with optimal growth occurring at pH 2–3 and temperatures of 75–80 °C, making them acidophiles and thermophiles respectively. Sulfolobus cells are irregularly shaped and flagellar.
Species of Sulfolobus are generally named after the location from which they were first isolated, e.g. Sulfolobus solfataricus was first isolated in the Solfatara volcano. Other species can be found throughout the world in areas of volcanic or geothermal activity, such as geological formations called mud pots, which are also known as solfatare (plural of solfatara).
Sulfolobus as a model to study the molecular mechanisms of DNA replication
When the first Archaeal genome, Methanococcus jannaschii, had been sequenced completely in 1996, it was found that the genes in the genome of Methanococcus jannaschii involved in DNA replication, transcription, and translation were more related to their counterparts in eukaryotes than to those in other prokaryotes. In 2001, the first genome sequence of Sulfolobus, Sulfolobus solfataricus P2, was published. In P2's genome, the genes related to chromosome replication were likewise found to be more related to those in eukaryotes. These genes include DNA polymerase, primase (including two subunits), MCM, CDC6/ORC1, RPA, RPC, and PCNA. In 2004, the origins of DNA replication of Sulfolobus solfataricus and Sulfolobus acidocaldarius were identified. It showed that both species contained two origins in their genome. This was the first time that more than a single origin of DNA replication had been shown to be used in a prokaryotic cell. The mechanism of DNA replication in archaea is evolutionary conserved, and similar to that of eukaryotes. Sulfolobus is now used as a model to study the molecular mechanisms of DNA replication in Archaea. And because the system of DNA replication in Archaea is much simpler than that in Eukaryota, it was suggested that Archaea could be used as a model to study the much more complex DNA replication in Eukaryota.
Role in biotechnology
Sulfolobus proteins are of interest for biotechnology and industrial use due to their thermostable nature. One application is the creation of artificial derivatives from S. acidocaldarius proteins, named affitins. Intracellular proteins are not necessarily stable at low pH however, as Sulfolobus species maintain a significant pH gradient across the outer membrane. Sulfolobales are metabolically dependent on sulfur: heterotrophic or autotrophic, their energy comes from the oxidation of sulfur and/or cellular respiration in which sulfur acts as the final electron acceptor. For example, S. tokodaii is known to oxidize hydrogen sulfide to sulfate intracellularly.
Phylogeny
The currently accepted taxonomy is based on the List of Prokaryotic names with Standing in Nomenclature (LPSN) and National Center for Biotechnology Information (NCBI)
Genome status
The complete genomes have been sequenced for S. acidocaldarius DSM 639 (2,225,959 nucleotides), S. solfataricus P2 (2,992,245 nucleotides), and S. tokodaii str. 7 (2,694,756 nucleotides).
Genome structure
The archaeon Sulfolobus solfataricus has a circular chromosome that consists of 2,992,245 bp. Another sequenced species, S. tokodaii has a circular chromosome as well but is slightly smaller with 2,694,756 bp. Both species lack the genes ftsZ and minD, which has been characteristic of sequenced Crenarchaeota. They also code for citrate synthase and two subunits of 2-oxoacid:ferredoxin oxidoreductase, which plays the same role as alpha-ketoglutarate dehydrogenase in the TCA (tricarboxylic/Krebs/citric acid) cycle. This indicates that Sulfolobus has a TCA cycle system similar to that found in mitochondria of eukaryotes. Other genes in the respiratory chain which partake in the production of ATP were not similar to what is found in eukaryotes. Cytochrome c is one such example that plays an important role in electron transfer to oxygen in eukaryotes. This was also found in A. pernix K1. Since this step is important for an aerobic microorganism like Sulfolobus, it probably uses a different molecule for the same function or has a different pathway.
Cell structure and metabolism
Sulfolobus can grow either lithoautotrophically by oxidizing sulfur, or chemoheterotrophically using sulfur to oxidize simple reduced carbon compounds. Heterotrophic growth has only been observed, however, in the presence of oxygen. The principle metabolic pathways are a glycolytic pathway, a pentose phosphate pathway, and the TCA cycle.
All Archaea have lipids with ether links between the head group and side chains, making the lipids more resistant to heat and acidity than bacterial and eukaryotic ester-linked lipids. The Sulfolobales are known for unusual tetraether lipids. In Sulfolobales, the ether-linked lipids are joined covalently across the "bilayer," making tetraethers. Technically, therefore, the tetraethers form a monolayer, not a bilayer. The tetraethers help Sulfolobus species survive extreme acid as well as high temperature.
Ecology
S. solfataricus has been found in different areas including Yellowstone National Park, Mount St. Helens, Iceland, Italy, and Russia to name a few. Sulfolobus is located almost wherever there is volcanic activity. They thrive in environments where the temperature is about 80 °C with a pH at about 3 and sulfur present. Another species, S. tokodaii, has been located in an acidic spa in Beppu Hot Springs, Kyushu, Japan. Sediments from ~90m below the seafloor on the Peruvian continental margin are dominated by intact archaeal tetraethers, and a significant fraction of the community is sedimentary archaea taxonomically linked to the crenarchaeal Sulfolobales (Sturt, et al., 2004).
DNA damage response
Exposure of Sulfolobus solfataricus or Sulfolobus acidocaldarius to the DNA damaging agents UV-irradiation, bleomycin or mitomycin C induced cellular aggregation. Other physical stressors, such as pH or temperature shift, did not induce aggregation, suggesting that induction of aggregation is caused specifically by DNA damage. Ajon et al. showed that UV-induced cellular aggregation mediates chromosomal marker exchange with high frequency in S. acidocaldarius. Recombination rates exceeded those of uninduced cultures by up to three orders of magnitude. Wood et al. also showed that UV-irradiation increased the frequency of recombination due to genetic exchange in S. acidocaldarius. Frols et al. and Ajon et al. hypothesized that the UV-inducible DNA transfer process and subsequent homologous recombinational repair represents an important mechanism to maintain chromosome integrity in S. acidocaldarius and S. solfataricus. This response may be a primitive form of sexual interaction, similar to the more well-studied bacterial transformation that is also associated with DNA transfer between cells leading to homologous recombinational repair of DNA damage.
The ups operon
The ups operon of Sulfolobus species is highly induced by UV irradiation. The pili encoded by this operon are employed in promoting cellular aggregation, which is necessary for subsequent DNA exchange between cells, resulting in homologous recombination. A study of the Sulfolobales acidocaldarius ups operon showed that one of the genes of the operon, saci-1497, encodes an endonuclease III that nicks UV-damaged DNA; and another gene of the operon, saci-1500, encodes a RecQ-like helicase that is able to unwind homologous recombination intermediates such as Holliday junctions. It was proposed that Saci-1497 and Saci-1500 function in an homologous recombination-based DNA repair mechanism that uses transferred DNA as a template. Thus it is thought that the ups system in combination with homologous recombination provide a DNA damage response which rescues Sulfolobales from DNA damaging threats.
Sulfolobus as a viral host
Lysogenic viruses infect Sulfolobus for protection. The viruses cannot survive in the extremely acidic and hot conditions that Sulfolobus lives in, and so the viruses use Sulfolobus as protection against the harsh elements. This relationship allows the virus to replicate inside the archaea without being destroyed by the environment. The Sulfolobus viruses are temperate or permanent lysogens. Permanent lysogens differ from lysogenic bacteriophages in that the host cells are not lysed after the induction of Fuselloviridae production and eventually return to the lysogenic state. They are also unique in the sense that the genes encoding the structural proteins of the virus are constantly transcribed and DNA replication appears to be induced. The viruses infecting archaea like Sulfolobus have to use a strategy to escape prolonged direct exposure to the type of environment their host lives in, which may explain some of their unique properties.
| Biology and health sciences | Archaea | Plants |
5839439 | https://en.wikipedia.org/wiki/Rotorcraft | Rotorcraft | A rotary-wing aircraft, rotorwing aircraft or rotorcraft is a heavier-than-air aircraft with rotary wings that spin around a vertical mast to generate lift. Part 1 (Definitions and Abbreviations) of Subchapter A of Chapter I of Title 14 of the U. S. Code of Federal Regulations states that rotorcraft "means a heavier-than-air aircraft that depends principally for its support in flight on the lift generated by one or more rotors." The assembly of several rotor blades mounted on a single mast is referred to as a rotor. The International Civil Aviation Organization (ICAO) defines a rotorcraft as "supported in flight by the reactions of the air on one or more rotors".
Rotorcraft generally include aircraft where one or more rotors provide lift throughout the entire flight, such as helicopters, gyroplanes, autogyros, and gyrodynes Compound rotorcraft augment the rotor with additional thrust engines, propellers, or static lifting surfaces. Some types, such as helicopters, are capable of vertical takeoff and landing. An aircraft which uses rotor lift for vertical flight but changes to solely fixed-wing lift in horizontal flight is not a rotorcraft but a convertiplane.
Classes of rotorcraft
Helicopter
A helicopter is a powered rotorcraft with rotors driven by the engine(s) throughout the flight, allowing it to take off and land vertically, hover, and fly forward, backward, or laterally. Helicopters have several different configurations of one or more main rotors.
Helicopters with a single shaft-driven main lift rotor require some sort of antitorque device such as a tail rotor, fantail, or NOTAR, except some rare examples of helicopters using tip jet propulsion, which generates almost no torque.
Autogyro
An autogyro (sometimes called gyrocopter, gyroplane, or rotaplane) uses an unpowered rotor, driven by aerodynamic forces in a state of autorotation to develop lift, and an engine-powered propeller, similar to that of a fixed-wing aircraft, to provide thrust. While similar to a helicopter rotor in appearance, the autogyro's rotor must have air flowing up and through the rotor disk in order to generate rotation. Early autogyros resembled the fixed-wing aircraft of the day, with wings and a front-mounted engine and propeller in a tractor configuration to pull the aircraft through the air. Late-model autogyros feature a rear-mounted engine and propeller in a pusher configuration.
The autogyro was invented in 1920 by Juan de la Cierva. The autogyro with pusher propeller was first tested by Etienne Dormoy with his Buhl A-1 Autogyro.
Gyrodyne
The rotor of a gyrodyne is normally driven by its engine for takeoff and landinghovering like a helicopterwith anti-torque and propulsion for forward flight provided by one or more propellers mounted on short or stub wings. As power is increased to the propeller, less power is required by the rotor to provide forward thrust resulting in reduced pitch angles and rotor blade flapping. At cruise speeds with most or all of the thrust being provided by the propellers, the rotor receives power only sufficient to overcome the profile drag and maintain lift. The effect is a rotorcraft operating in a more efficient manner than the freewheeling rotor of an autogyro in autorotation, minimizing the adverse effects of retreating blade stall of helicopters at higher airspeeds.
Rotor kite
A rotor kite or gyroglider is an unpowered rotary-wing aircraft. Like an autogyro or helicopter, it relies on lift created by one or more sets of rotors in order to fly. Unlike a helicopter, autogyros and rotor kites do not have an engine powering their rotors, but while an autogyro has an engine providing forward thrust that keeps the rotor turning, a rotor kite has no engine at all, and relies on either being carried aloft and dropped from another aircraft, or by being towed into the air behind a car or boat.
Rotor configuration
Number of blades
A rotary wing is characterised by the number of blades. Typically this is between two and six per driveshaft.
Number of rotors
A rotorcraft may have one or more rotors. Various rotor configurations have been used:
One rotor. Powered rotors require compensation for the torque reaction causing yaw, except in the case of tipjet drive. One rotor rotorcraft are typically called monocopters.
Two rotors. These typically rotate in opposite directions cancelling the torque reaction so that no tail rotor or other yaw stabiliser is needed. These rotors can be laid out as
TandemOne in front of the other.
TransverseSide by side.
CoaxialOne rotor disc above the other, with concentric drive shafts.
IntermeshingTwin rotors at an acute angle from each other, whose nearly-vertical driveshafts are geared together to synchronise their rotor blades so that they intermesh, also called a synchropter.
Three rotors. An uncommon configuration; the 1948 Cierva Air Horse had three rotors as it was not believed a single rotor of sufficient strength could be built for its size. All three rotors turned in the same direction and yaw compensation was provided by inclining each rotor axis to generate rotor thrust components that opposed torque.
Four rotors. Also referred to as the quadcopter or quadrotor. Usually two rotors turn clockwise and two counter-clockwise.
More than four rotors. Referred to generally as multirotors, or sometimes individually as hexacopters and octocopter, these configurations typically have matched sets of rotors turning in opposite directions. They are uncommon in full-size manned aircraft but are popular for unmanned aerial vehicles (UAVs).
Stopped rotors
Some rotary wing aircraft are designed to stop the rotor for forward flight so that it then acts as a fixed wing. For vertical flight and hovering it spins to act as a rotary wing or rotor, and for forward flight at speed it stops to act as a fixed wing providing some or all of the lift required. Additional fixed wings may also be provided to help with stability and control and to provide auxiliary lift.
An early American proposal was the conversion of the Lockheed F-104 Starfighter with a triangular rotor wing. The idea was later revisited by Hughes. The Sikorsky S-72 research aircraft underwent extensive flight testing.
In 1986 the Sikorsky S-72 Rotor Systems Research Aircraft (RSRA) was fitted with a four-bladed stopped rotor, known as the X-wing. The programme was cancelled two years later, before the rotor had flown.
The later canard rotor/wing (CRW) concept added a "canard" foreplane as well as a conventional tailplane, offloading the rotor wing and providing control during forward flight. For vertical and low-speed flight, the main airfoil is tip-driven as a helicopter's rotor by exhaust from a jet engine, and there is no need for a tail rotor. In high-speed flight the airfoil is stopped in a spanwise position, as the main wing of a three-surface aircraft, and the engine exhausts through an ordinary jet nozzle. Two Boeing X-50 Dragonfly prototypes with a two-bladed rotor were flown from 2003 but the program ended after both had crashed, having failed to transition successfully.
In 2013 the US Naval Research Laboratory (NRL) published a vertical-to-horizontal flight transition method and associated technology, patented December 6, 2011, which they call the Stop-Rotor Rotary Wing Aircraft. The Australian company StopRotor Technology Pty Ltd has developed a prototype Hybrid RotorWing (HRW) craft. The design uses high alpha airflow to provide a symmetrical airflow across all the rotor blades, requiring it to drop almost vertically during transition. Inflight transition from fixed to rotary mode was demonstrated in August 2013.
Another approach proposes a tailsitter configuration in which the lifting surfaces act as a rotors during takeoff, the craft tilts over for horizontal flight and the rotor stops to act as a fixed wing.
| Technology | Types of aircraft | null |
5839673 | https://en.wikipedia.org/wiki/Quantitative%20analysis%20%28chemistry%29 | Quantitative analysis (chemistry) | In analytical chemistry, quantitative analysis is the determination of the absolute or relative abundance (often expressed as a concentration) of one, several or all particular substance(s) present in a sample. It relates to the determination of percentage of constituents in any given sample.
Methods
Once the presence of certain substances in a sample is known, the study of their absolute or relative abundance could help in determining specific properties. Knowing the composition of a sample is very important, and several ways have been developed to make it possible, like gravimetric and volumetric analysis. Gravimetric analysis yields more accurate data about the composition of a sample than volumetric analysis but also takes more time to perform in the laboratory. Volumetric analysis, on the other hand, doesn't take that much time and can produce satisfactory results. Volumetric analysis can be simply a titration based in a neutralization reaction but it can also be a precipitation or a complex forming reaction as well as a titration based in a redox reaction. However, each method in quantitative analysis has a general specification, in neutralization reactions, for example, the reaction that occurs is between an acid and a base, which yields a salt and water, hence the name neutralization. In the precipitation reactions the standard solution is in the most cases silver nitrate which is used as a reagent to react with the ions present in the sample and to form a highly insoluble precipitate. Precipitation methods are often called simply as argentometry. In the two other methods the situation is the same. Complex forming titration is a reaction that occurs between metal ions and a standard solution that is in the most cases EDTA (Ethylene Diamine Tetra Acetic acid). In the redox titration that reaction is carried out between an oxidizing agent and a reduction agent. There are some more methods like Liebig method / Duma's method / Kjeldahl's method and Carius method for estimation of organic compounds.
For example, quantitative analysis performed by mass spectrometry on biological samples can determine, by the relative abundance ratio of specific proteins, indications of certain diseases, like cancer.
Quantitative vs. qualitative
The term "quantitative analysis" is often used in comparison (or contrast) with "qualitative analysis", which seeks information about the identity or form of substance present. For instance, a chemist might be given an unknown solid sample. They will use "qualitative" techniques (perhaps NMR or IR spectroscopy) to identify the compounds present, and then quantitative techniques to determine the amount of each compound in the sample. Careful procedures for recognizing the presence of different metal ions have been developed, although they have largely been replaced by modern instruments; these are collectively known as qualitative inorganic analysis. Similar tests for identifying organic compounds (by testing for different functional groups) are also known.
Many techniques can be used for either qualitative or quantitative measurements. For instance, suppose an indicator solution changes color in the presence of a metal ion. It could be used as a qualitative test: does the indicator solution change color when a drop of sample is added? It could also be used as a quantitative test, by studying the color of the indicator solution with different concentrations of the metal ion. (This would probably be done using ultraviolet-visible spectroscopy.)
| Physical sciences | Basics_2 | Chemistry |
390214 | https://en.wikipedia.org/wiki/Graphics%20processing%20unit | Graphics processing unit | A graphics processing unit (GPU) is a specialized electronic circuit initially designed for digital image processing and to accelerate computer graphics, being present either as a discrete video card or embedded on motherboards, mobile phones, personal computers, workstations, and game consoles. After their initial design, GPUs were found to be useful for non-graphic calculations involving embarrassingly parallel problems due to their parallel structure. Other non-graphical uses include the training of neural networks and cryptocurrency mining.
History
1970s
Arcade system boards have used specialized graphics circuits since the 1970s. In early video game hardware, RAM for frame buffers was expensive, so video chips composited data together as the display was being scanned out on the monitor.
A specialized barrel shifter circuit helped the CPU animate the framebuffer graphics for various 1970s arcade video games from Midway and Taito, such as Gun Fight (1975), Sea Wolf (1976), and Space Invaders (1978). The Namco Galaxian arcade system in 1979 used specialized graphics hardware that supported RGB color, multi-colored sprites, and tilemap backgrounds. The Galaxian hardware was widely used during the golden age of arcade video games, by game companies such as Namco, Centuri, Gremlin, Irem, Konami, Midway, Nichibutsu, Sega, and Taito.
The Atari 2600 in 1977 used a video shifter called the Television Interface Adaptor. Atari 8-bit computers (1979) had ANTIC, a video processor which interpreted instructions describing a "display list"—the way the scan lines map to specific bitmapped or character modes and where the memory is stored (so there did not need to be a contiguous frame buffer). 6502 machine code subroutines could be triggered on scan lines by setting a bit on a display list instruction. ANTIC also supported smooth vertical and horizontal scrolling independent of the CPU.
1980s
The NEC μPD7220 was the first implementation of a personal computer graphics display processor as a single large-scale integration (LSI) integrated circuit chip. This enabled the design of low-cost, high-performance video graphics cards such as those from Number Nine Visual Technology. It became the best-known GPU until the mid-1980s. It was the first fully integrated VLSI (very large-scale integration) metal–oxide–semiconductor (NMOS) graphics display processor for PCs, supported up to 1024×1024 resolution, and laid the foundations for the emerging PC graphics market. It was used in a number of graphics cards and was licensed for clones such as the Intel 82720, the first of Intel's graphics processing units. The Williams Electronics arcade games Robotron 2084, Joust, Sinistar, and Bubbles, all released in 1982, contain custom blitter chips for operating on 16-color bitmaps.
In 1984, Hitachi released ARTC HD63484, the first major CMOS graphics processor for personal computers. The ARTC could display up to 4K resolution when in monochrome mode. It was used in a number of graphics cards and terminals during the late 1980s. In 1985, the Amiga was released with a custom graphics chip including a blitter for bitmap manipulation, line drawing, and area fill. It also included a coprocessor with its own simple instruction set, that was capable of manipulating graphics hardware registers in sync with the video beam (e.g. for per-scanline palette switches, sprite multiplexing, and hardware windowing), or driving the blitter. In 1986, Texas Instruments released the TMS34010, the first fully programmable graphics processor. It could run general-purpose code, but it had a graphics-oriented instruction set. During 1990–1992, this chip became the basis of the Texas Instruments Graphics Architecture ("TIGA") Windows accelerator cards.
In 1987, the IBM 8514 graphics system was released. It was one of the first video cards for IBM PC compatibles to implement fixed-function 2D primitives in electronic hardware. Sharp's X68000, released in 1987, used a custom graphics chipset with a 65,536 color palette and hardware support for sprites, scrolling, and multiple playfields. It served as a development machine for Capcom's CP System arcade board. Fujitsu's FM Towns computer, released in 1989, had support for a 16,777,216 color palette. In 1988, the first dedicated polygonal 3D graphics boards were introduced in arcades with the Namco System 21 and Taito Air System.
IBM introduced its proprietary Video Graphics Array (VGA) display standard in 1987, with a maximum resolution of 640×480 pixels. In November 1988, NEC Home Electronics announced its creation of the Video Electronics Standards Association (VESA) to develop and promote a Super VGA (SVGA) computer display standard as a successor to VGA. Super VGA enabled graphics display resolutions up to 800×600 pixels, a 36% increase.
1990s
In 1991, S3 Graphics introduced the S3 86C911, which its designers named after the Porsche 911 as an indication of the performance increase it promised. The 86C911 spawned a variety of imitators: by 1995, all major PC graphics chip makers had added 2D acceleration support to their chips. Fixed-function Windows accelerators surpassed expensive general-purpose graphics coprocessors in Windows performance, and such coprocessors faded from the PC market.
Throughout the 1990s, 2D GUI acceleration evolved. As manufacturing capabilities improved, so did the level of integration of graphics chips. Additional application programming interfaces (APIs) arrived for a variety of tasks, such as Microsoft's WinG graphics library for Windows 3.x, and their later DirectDraw interface for hardware acceleration of 2D games in Windows 95 and later.
In the early- and mid-1990s, real-time 3D graphics became increasingly common in arcade, computer, and console games, which led to increasing public demand for hardware-accelerated 3D graphics. Early examples of mass-market 3D graphics hardware can be found in arcade system boards such as the Sega Model 1, Namco System 22, and Sega Model 2, and the fifth-generation video game consoles such as the Saturn, PlayStation, and Nintendo 64. Arcade systems such as the Sega Model 2 and SGI Onyx-based Namco Magic Edge Hornet Simulator in 1993 were capable of hardware T&L (transform, clipping, and lighting) years before appearing in consumer graphics cards. Another early example is the Super FX chip, a RISC-based on-cartridge graphics chip used in some SNES games, notably Doom and Star Fox. Some systems used DSPs to accelerate transformations. Fujitsu, which worked on the Sega Model 2 arcade system, began working on integrating T&L into a single LSI solution for use in home computers in 1995; the Fujitsu Pinolite, the first 3D geometry processor for personal computers, released in 1997. The first hardware T&L GPU on home video game consoles was the Nintendo 64's Reality Coprocessor, released in 1996. In 1997, Mitsubishi released the 3Dpro/2MP, a GPU capable of transformation and lighting, for workstations and Windows NT desktops; ATi used it for its FireGL 4000 graphics card, released in 1997.
The term "GPU" was coined by Sony in reference to the 32-bit Sony GPU (designed by Toshiba) in the PlayStation video game console, released in 1994.
In the PC world, notable failed attempts for low-cost 3D graphics chips included the S3 ViRGE, ATI Rage, and Matrox Mystique. These chips were essentially previous-generation 2D accelerators with 3D features bolted on. Many were pin-compatible with the earlier-generation chips for ease of implementation and minimal cost. Initially, 3D graphics were possible only with discrete boards dedicated to accelerating 3D functions (and lacking 2D graphical user interface (GUI) acceleration entirely) such as the PowerVR and the 3dfx Voodoo. However, as manufacturing technology continued to progress, video, 2D GUI acceleration, and 3D functionality were all integrated into one chip. Rendition's Verite chipsets were among the first to do this well. In 1997, Rendition collaborated with Hercules and Fujitsu on a "Thriller Conspiracy" project which combined a Fujitsu FXG-1 Pinolite geometry processor with a Vérité V2200 core to create a graphics card with a full T&L engine years before Nvidia's GeForce 256; This card, designed to reduce the load placed upon the system's CPU, never made it to market. NVIDIA RIVA 128 was one of the first consumer-facing GPU integrated 3D processing unit and 2D processing unit on a chip.
OpenGL was introduced in the early '90s by SGI as a professional graphics API, with proprietary hardware support for 3D rasterization. In 1994 Microsoft acquired Softimage, the dominant CGI movie production tool used for early CGI movie hits like Jurassic Park, Terminator 2 and Titanic. With that deal came a strategic relationship with SGI and a commercial license of SGI's OpenGL libraries enabling Microsoft to port the API to the Windows NT OS but not to the upcoming release of Windows '95. Although it was little known at the time, SGI had contracted with Microsoft to transition from Unix to the forthcoming Windows NT OS, the deal which was signed in 1995 was not announced publicly until 1998. In the intervening period, Microsoft worked closely with SGI to port OpenGL to Windows NT. In that era OpenGL had no standard driver model for competing hardware accelerators to compete on the basis of support for higher level 3D texturing and lighting functionality. In 1994 Microsoft announced DirectX 1.0 and support for gaming in the forthcoming Windows '95 consumer OS, in '95 Microsoft announced the acquisition of UK based Rendermorphics Ltd and the Direct3D driver model for the acceleration of consumer 3D graphics. The Direct3D driver model shipped with DirectX 2.0 in 1996. It included standards and specifications for 3D chip makers to compete to support 3D texture, lighting and Z-buffering. ATI, which was later to be acquired by AMD, began development on the first Direct3D GPU's. Nvidia, quickly pivoted from a failed deal with Sega in 1996 to aggressively embracing support for Direct3D. In this era Microsoft merged their internal Direct3D and OpenGL teams and worked closely with SGI to unify driver standards for both industrial and consumer 3D graphics hardware accelerators. Microsoft ran annual events for 3D chip makers called "Meltdowns" to test their 3D hardware and drivers to work both with Direct3D and OpenGL. It was during this period of strong Microsoft influence over 3D standards that 3D accelerator cards moved beyond being simple rasterizers to become more powerful general purpose processors as support for hardware accelerated texture mapping, lighting, Z-buffering and compute created the modern GPU. During this period the same Microsoft team responsible for Direct3D and OpenGL driver standardization introduced their own Microsoft 3D chip design called Talisman. Details of this era are documented extensively in the books: "Game of X" v.1 and v.2 by Russel Demaria, "Renegades of the Empire" by Mike Drummond, "Opening the Xbox" by Dean Takahashi and "Masters of Doom" by David Kushner. The Nvidia GeForce 256 (also known as NV10) was the first consumer-level card with hardware-accelerated T&L; While the OpenGL API provided software support for texture mapping and lighting the first 3D hardware acceleration for these features arrived with the first Direct3D accelerated consumer GPU's.
2000s
Nvidia was first to produce a chip capable of programmable shading: the GeForce 3. Each pixel could now be processed by a short program that could include additional image textures as inputs, and each geometric vertex could likewise be processed by a short program before it was projected onto the screen. Used in the Xbox console, this chip competed with the one in the PlayStation 2, which used a custom vector unit for hardware accelerated vertex processing (commonly referred to as VU0/VU1). The earliest incarnations of shader execution engines used in Xbox were not general purpose and could not execute arbitrary pixel code. Vertices and pixels were processed by different units which had their own resources, with pixel shaders having tighter constraints (because they execute at higher frequencies than vertices). Pixel shading engines were actually more akin to a highly customizable function block and did not really "run" a program. Many of these disparities between vertex and pixel shading were not addressed until the Unified Shader Model.
In October 2002, with the introduction of the ATI Radeon 9700 (also known as R300), the world's first Direct3D 9.0 accelerator, pixel and vertex shaders could implement looping and lengthy floating point math, and were quickly becoming as flexible as CPUs, yet orders of magnitude faster for image-array operations. Pixel shading is often used for bump mapping, which adds texture to make an object look shiny, dull, rough, or even round or extruded.
With the introduction of the Nvidia GeForce 8 series and new generic stream processing units, GPUs became more generalized computing devices. Parallel GPUs are making computational inroads against the CPU, and a subfield of research, dubbed GPU computing or GPGPU for general purpose computing on GPU, has found applications in fields as diverse as machine learning, oil exploration, scientific image processing, linear algebra, statistics, 3D reconstruction, and stock options pricing. GPGPU was the precursor to what is now called a compute shader (e.g. CUDA, OpenCL, DirectCompute) and actually abused the hardware to a degree by treating the data passed to algorithms as texture maps and executing algorithms by drawing a triangle or quad with an appropriate pixel shader. This entails some overheads since units like the scan converter are involved where they are not needed (nor are triangle manipulations even a concern—except to invoke the pixel shader).
Nvidia's CUDA platform, first introduced in 2007, was the earliest widely adopted programming model for GPU computing. OpenCL is an open standard defined by the Khronos Group that allows for the development of code for both GPUs and CPUs with an emphasis on portability. OpenCL solutions are supported by Intel, AMD, Nvidia, and ARM, and according to a report in 2011 by Evans Data, OpenCL had become the second most popular HPC tool.
2010s
In 2010, Nvidia partnered with Audi to power their cars' dashboards, using the Tegra GPU to provide increased functionality to cars' navigation and entertainment systems. Advances in GPU technology in cars helped advance self-driving technology. AMD's Radeon HD 6000 series cards were released in 2010, and in 2011 AMD released its 6000M Series discrete GPUs for mobile devices. The Kepler line of graphics cards by Nvidia were released in 2012 and were used in the Nvidia's 600 and 700 series cards. A feature in this GPU microarchitecture included GPU boost, a technology that adjusts the clock-speed of a video card to increase or decrease it according to its power draw. The Kepler microarchitecture was manufactured .
The PS4 and Xbox One were released in 2013; they both use GPUs based on AMD's Radeon HD 7850 and 7790. Nvidia's Kepler line of GPUs was followed by the Maxwell line, manufactured on the same process. Nvidia's 28 nm chips were manufactured by TSMC in Taiwan using the 28 nm process. Compared to the 40 nm technology from the past, this manufacturing process allowed a 20 percent boost in performance while drawing less power. Virtual reality headsets have high system requirements; manufacturers recommended the GTX 970 and the R9 290X or better at the time of their release. Cards based on the Pascal microarchitecture were released in 2016. The GeForce 10 series of cards are of this generation of graphics cards. They are made using the 16 nm manufacturing process which improves upon previous microarchitectures. Nvidia released one non-consumer card under the new Volta architecture, the Titan V. Changes from the Titan XP, Pascal's high-end card, include an increase in the number of CUDA cores, the addition of tensor cores, and HBM2. Tensor cores are designed for deep learning, while high-bandwidth memory is on-die, stacked, lower-clocked memory that offers an extremely wide memory bus. To emphasize that the Titan V is not a gaming card, Nvidia removed the "GeForce GTX" suffix it adds to consumer gaming cards.
In 2018, Nvidia launched the RTX 20 series GPUs that added ray-tracing cores to GPUs, improving their performance on lighting effects. Polaris 11 and Polaris 10 GPUs from AMD are fabricated by a 14 nm process. Their release resulted in a substantial increase in the performance per watt of AMD video cards. AMD also released the Vega GPU series for the high end market as a competitor to Nvidia's high end Pascal cards, also featuring HBM2 like the Titan V.
In 2019, AMD released the successor to their Graphics Core Next (GCN) microarchitecture/instruction set. Dubbed RDNA, the first product featuring it was the Radeon RX 5000 series of video cards.
The company announced that the successor to the RDNA microarchitecture would be incremental (aka a refresh). AMD unveiled the Radeon RX 6000 series, its RDNA 2 graphics cards with support for hardware-accelerated ray tracing. The product series, launched in late 2020, consisted of the RX 6800, RX 6800 XT, and RX 6900 XT. The RX 6700 XT, which is based on Navi 22, was launched in early 2021.
The PlayStation 5 and Xbox Series X and Series S were released in 2020; they both use GPUs based on the RDNA 2 microarchitecture with incremental improvements and different GPU configurations in each system's implementation.
Intel first entered the GPU market in the late 1990s, but produced lackluster 3D accelerators compared to the competition at the time. Rather than attempting to compete with the high-end manufacturers Nvidia and ATI/AMD, they began integrating Intel Graphics Technology GPUs into motherboard chipsets, beginning with the Intel 810 for the Pentium III, and later into CPUs. They began with the Intel Atom 'Pineview' laptop processor in 2009, continuing in 2010 with desktop processors in the first generation of the Intel Core line and with contemporary Pentiums and Celerons. This resulted in a large nominal market share, as the majority of computers with an Intel CPU also featured this embedded graphics processor. These generally lagged behind discrete processors in performance. Intel re-entered the discrete GPU market in 2022 with its Arc series, which competed with the then-current GeForce 30 series and Radeon 6000 series cards at competitive prices.
2020s
In the 2020s, GPUs have been increasingly used for calculations involving embarrassingly parallel problems, such as training of neural networks on enormous datasets that are needed for large language models. Specialized processing cores on some modern workstation's GPUs are dedicated for deep learning since they have significant FLOPS performance increases, using 4×4 matrix multiplication and division, resulting in hardware performance up to 128 TFLOPS in some applications. These tensor cores are expected to appear in consumer cards, as well.
GPU companies
Many companies have produced GPUs under a number of brand names. In 2009, Intel, Nvidia, and AMD/ATI were the market share leaders, with 49.4%, 27.8%, and 20.6% market share respectively. In addition, Matrox produces GPUs. Chinese companies such as Jingjia Micro have also produced GPUs for the domestic market although in terms of worldwide sales, they still lag behind market leaders.
Modern smartphones use mostly Adreno GPUs from Qualcomm, PowerVR GPUs from Imagination Technologies, and Mali GPUs from ARM.
Computational functions
Modern GPUs have traditionally used most of their transistors to do calculations related to 3D computer graphics. In addition to the 3D hardware, today's GPUs include basic 2D acceleration and framebuffer capabilities (usually with a VGA compatibility mode). Newer cards such as AMD/ATI HD5000–HD7000 lack dedicated 2D acceleration; it is emulated by 3D hardware. GPUs were initially used to accelerate the memory-intensive work of texture mapping and rendering polygons. Later, were added to accelerate geometric calculations such as the rotation and translation of vertices into different coordinate systems. Recent developments in GPUs include support for programmable shaders which can manipulate vertices and textures with many of the same operations that are supported by CPUs, oversampling and interpolation techniques to reduce aliasing, and very high-precision color spaces.
Several factors of GPU construction affect the performance of the card for real-time rendering, such as the size of the connector pathways in the semiconductor device fabrication, the clock signal frequency, and the number and size of various on-chip memory caches. Performance is also affected by the number of streaming multiprocessors (SM) for NVidia GPUs, or compute units (CU) for AMD GPUs, or Xe cores for Intel discrete GPUs, which describe the number of on-silicon processor core units within the GPU chip that perform the core calculations, typically working in parallel with other SM/CUs on the GPU. GPU performance is typically measured in floating point operations per second (FLOPS); GPUs in the 2010s and 2020s typically deliver performance measured in teraflops (TFLOPS). This is an estimated performance measure, as other factors can affect the actual display rate.
GPU accelerated video decoding and encoding
Most GPUs made since 1995 support the YUV color space and hardware overlays, important for digital video playback, and many GPUs made since 2000 also support MPEG primitives such as motion compensation and iDCT. This hardware-accelerated video decoding, in which portions of the video decoding process and video post-processing are offloaded to the GPU hardware, is commonly referred to as "GPU accelerated video decoding", "GPU assisted video decoding", "GPU hardware accelerated video decoding", or "GPU hardware assisted video decoding".
Recent graphics cards decode high-definition video on the card, offloading the central processing unit. The most common APIs for GPU accelerated video decoding are DxVA for Microsoft Windows operating systems and VDPAU, VAAPI, XvMC, and XvBA for Linux-based and UNIX-like operating systems. All except XvMC are capable of decoding videos encoded with MPEG-1, MPEG-2, MPEG-4 ASP (MPEG-4 Part 2), MPEG-4 AVC (H.264 / DivX 6), VC-1, WMV3/WMV9, Xvid / OpenDivX (DivX 4), and DivX 5 codecs, while XvMC is only capable of decoding MPEG-1 and MPEG-2.
There are several dedicated hardware video decoding and encoding solutions.
Video decoding processes that can be accelerated
Video decoding processes that can be accelerated by modern GPU hardware are:
Motion compensation (mocomp)
Inverse discrete cosine transform (iDCT)
Inverse telecine 3:2 and 2:2 pull-down correction
Inverse modified discrete cosine transform (iMDCT)
In-loop deblocking filter
Intra-frame prediction
Inverse quantization (IQ)
Variable-length decoding (VLD), more commonly known as slice-level acceleration
Spatial-temporal deinterlacing and automatic interlace/progressive source detection
Bitstream processing (Context-adaptive variable-length coding/Context-adaptive binary arithmetic coding) and perfect pixel positioning
These operations also have applications in video editing, encoding, and transcoding.
2D graphics APIs
An earlier GPU may support one or more 2D graphics API for 2D acceleration, such as GDI and DirectDraw.
3D graphics APIs
A GPU can support one or more 3D graphics API, such as DirectX, Metal, OpenGL, OpenGL ES, Vulkan.
GPU forms
Terminology
In the 1970s, the term "GPU" originally stood for graphics processor unit and described a programmable processing unit working independently from the CPU that was responsible for graphics manipulation and output. In 1994, Sony used the term (now standing for graphics processing unit) in reference to the PlayStation console's Toshiba-designed Sony GPU. The term was popularized by Nvidia in 1999, who marketed the GeForce 256 as "the world's first GPU". It was presented as a "single-chip processor with integrated transform, lighting, triangle setup/clipping, and rendering engines". Rival ATI Technologies coined the term "visual processing unit" or VPU with the release of the Radeon 9700 in 2002. The AMD Alveo MA35D features dual VPU’s, each using the 5 nm process in 2023.
In personal computers, there are two main forms of GPUs. Each has many synonyms:
Dedicated graphics also called discrete graphics.
Integrated graphics also called shared graphics solutions, integrated graphics processors (IGP), or unified memory architecture (UMA).
Usage-specific GPU
Most GPUs are designed for a specific use, real-time 3D graphics, or other mass calculations:
Gaming
GeForce GTX, RTX
Nvidia Titan
Radeon HD, R5, R7, R9, RX, Vega and Navi series
Radeon VII
Intel Arc
Cloud Gaming
Nvidia GRID
Radeon Sky
Workstation
Nvidia Quadro
Nvidia RTX
AMD FirePro
AMD Radeon Pro
Intel Arc Pro
Cloud Workstation
Nvidia Tesla
AMD FireStream
Artificial Intelligence training and Cloud
Nvidia Tesla
AMD Radeon Instinct
Automated/Driverless car
Nvidia Drive PX
Dedicated graphics processing unit
Dedicated graphics processing units uses RAM that is dedicated to the GPU rather than relying on the computer’s main system memory. This RAM is usually specially selected for the expected serial workload of the graphics card (see GDDR). Sometimes systems with dedicated discrete GPUs were called "DIS" systems as opposed to "UMA" systems (see next section).
Dedicated GPUs are not necessarily removable, nor does it necessarily interface with the motherboard in a standard fashion. The term "dedicated" refers to the fact that graphics cards have RAM that is dedicated to the card's use, not to the fact that most dedicated GPUs are removable. Dedicated GPUs for portable computers are most commonly interfaced through a non-standard and often proprietary slot due to size and weight constraints. Such ports may still be considered PCIe or AGP in terms of their logical host interface, even if they are not physically interchangeable with their counterparts.
Graphics cards with dedicated GPUs typically interface with the motherboard by means of an expansion slot such as PCI Express (PCIe) or Accelerated Graphics Port (AGP). They can usually be replaced or upgraded with relative ease, assuming the motherboard is capable of supporting the upgrade. A few graphics cards still use Peripheral Component Interconnect (PCI) slots, but their bandwidth is so limited that they are generally used only when a PCIe or AGP slot is not available.
Technologies such as Scan-Line Interleave by 3dfx, SLI and NVLink by Nvidia and CrossFire by AMD allow multiple GPUs to draw images simultaneously for a single screen, increasing the processing power available for graphics. These technologies, however, are increasingly uncommon; most games do not fully use multiple GPUs, as most users cannot afford them. Multiple GPUs are still used on supercomputers (like in Summit), on workstations to accelerate video (processing multiple videos at once) and 3D rendering, for VFX, GPGPU workloads and for simulations, and in AI to expedite training, as is the case with Nvidia's lineup of DGX workstations and servers, Tesla GPUs, and Intel's Ponte Vecchio GPUs.
Integrated graphics processing unit
Integrated graphics processing units (IGPU), integrated graphics, shared graphics solutions, integrated graphics processors (IGP), or unified memory architectures (UMA) use a portion of a computer's system RAM rather than dedicated graphics memory. IGPs can be integrated onto a motherboard as part of its northbridge chipset, or on the same die (integrated circuit) with the CPU (like AMD APU or Intel HD Graphics). On certain motherboards, AMD's IGPs can use dedicated sideport memory: a separate fixed block of high performance memory that is dedicated for use by the GPU. computers with integrated graphics account for about 90% of all PC shipments. They are less costly to implement than dedicated graphics processing, but tend to be less capable. Historically, integrated processing was considered unfit for 3D games or graphically intensive programs but could run less intensive programs such as Adobe Flash. Examples of such IGPs would be offerings from SiS and VIA circa 2004. However, modern integrated graphics processors such as AMD Accelerated Processing Unit and Intel Graphics Technology (HD, UHD, Iris, Iris Pro, Iris Plus, and Xe-LP) can handle 2D graphics or low-stress 3D graphics.
Since GPU computations are memory-intensive, integrated processing may compete with the CPU for relatively slow system RAM, as it has minimal or no dedicated video memory. IGPs use system memory with bandwidth up to a current maximum of 128 GB/s, whereas a discrete graphics card may have a bandwidth of more than 1000 GB/s between its VRAM and GPU core. This memory bus bandwidth can limit the performance of the GPU, though multi-channel memory can mitigate this deficiency. Older integrated graphics chipsets lacked hardware transform and lighting, but newer ones include it.
On systems with "Unified Memory Architecture" (UMA), including modern AMD processors with integrated graphics, modern Intel processors with integrated graphics, Apple processors, the PS5 and Xbox Series (among others), the CPU cores and the GPU block share the same pool of RAM and memory address space. This allows the system to dynamically allocate memory between the CPU cores and the GPU block based on memory needs (without needing a large static split of the RAM) and thanks to zero copy transfers, removes the need for either copying data over a bus (computing) between physically separate RAM pools or copying between separate address spaces on a single physical pool of RAM, allowing more efficient transfer of data.
Hybrid graphics processing
Hybrid GPUs compete with integrated graphics in the low-end desktop and notebook markets. The most common implementations of this are ATI's HyperMemory and Nvidia's TurboCache.
Hybrid graphics cards are somewhat more expensive than integrated graphics, but much less expensive than dedicated graphics cards. They share memory with the system and have a small dedicated memory cache, to make up for the high latency of the system RAM. Technologies within PCI Express make this possible. While these solutions are sometimes advertised as having as much as 768 MB of RAM, this refers to how much can be shared with the system memory.
Stream processing and general purpose GPUs (GPGPU)
It is common to use a general purpose graphics processing unit (GPGPU) as a modified form of stream processor (or a vector processor), running compute kernels. This turns the massive computational power of a modern graphics accelerator's shader pipeline into general-purpose computing power. In certain applications requiring massive vector operations, this can yield several orders of magnitude higher performance than a conventional CPU. The two largest discrete (see "Dedicated graphics processing unit" above) GPU designers, AMD and Nvidia, are pursuing this approach with an array of applications. Both Nvidia and AMD teamed with Stanford University to create a GPU-based client for the Folding@home distributed computing project for protein folding calculations. In certain circumstances, the GPU calculates forty times faster than the CPUs traditionally used by such applications.
GPGPUs can be used for many types of embarrassingly parallel tasks including ray tracing. They are generally suited to high-throughput computations that exhibit data-parallelism to exploit the wide vector width SIMD architecture of the GPU.
GPU-based high performance computers play a significant role in large-scale modelling. Three of the ten most powerful supercomputers in the world take advantage of GPU acceleration.
GPUs support API extensions to the C programming language such as OpenCL and OpenMP. Furthermore, each GPU vendor introduced its own API which only works with their cards: AMD APP SDK from AMD, and CUDA from Nvidia. These allow functions called compute kernels to run on the GPU's stream processors. This makes it possible for C programs to take advantage of a GPU's ability to operate on large buffers in parallel, while still using the CPU when appropriate. CUDA was the first API to allow CPU-based applications to directly access the resources of a GPU for more general purpose computing without the limitations of using a graphics API.
Since 2005 there has been interest in using the performance offered by GPUs for evolutionary computation in general, and for accelerating the fitness evaluation in genetic programming in particular. Most approaches compile linear or tree programs on the host PC and transfer the executable to the GPU to be run. Typically a performance advantage is only obtained by running the single active program simultaneously on many example problems in parallel, using the GPU's SIMD architecture. However, substantial acceleration can also be obtained by not compiling the programs, and instead transferring them to the GPU, to be interpreted there. Acceleration can then be obtained by either interpreting multiple programs simultaneously, simultaneously running multiple example problems, or combinations of both. A modern GPU can simultaneously interpret hundreds of thousands of very small programs.
External GPU (eGPU)
An external GPU is a graphics processor located outside of the housing of the computer, similar to a large external hard drive. External graphics processors are sometimes used with laptop computers. Laptops might have a substantial amount of RAM and a sufficiently powerful central processing unit (CPU), but often lack a powerful graphics processor, and instead have a less powerful but more energy-efficient on-board graphics chip. On-board graphics chips are often not powerful enough for playing video games, or for other graphically intensive tasks, such as editing video or 3D animation/rendering.
Therefore, it is desirable to attach a GPU to some external bus of a notebook. PCI Express is the only bus used for this purpose. The port may be, for example, an ExpressCard or mPCIe port (PCIe ×1, up to 5 or 2.5 Gbit/s respectively), a Thunderbolt 1, 2, or 3 port (PCIe ×4, up to 10, 20, or 40 Gbit/s respectively), a USB4 port with Thunderbolt compatibility, or an OCuLink port. Those ports are only available on certain notebook systems. eGPU enclosures include their own power supply (PSU), because powerful GPUs can consume hundreds of watts.
Official vendor support for external GPUs has gained traction. A milestone was Apple's decision to support external GPUs with MacOS High Sierra 10.13.4. Several major hardware vendors (HP, Razer) released Thunderbolt 3 eGPU enclosures. This support fuels eGPU implementations by enthusiasts.
Energy efficiency
Sales
In 2013, 438.3 million GPUs were shipped globally and the forecast for 2014 was 414.2 million. However, by the third quarter of 2022, shipments of integrated GPUs totaled around 75.5 million units, down 19% year-over-year.
| Technology | Computer hardware | null |
390410 | https://en.wikipedia.org/wiki/Salt%20pan%20%28geology%29 | Salt pan (geology) | Natural salt pans or salt flats are flat expanses of ground covered with salt and other minerals, usually shining white under the sun. They are found in deserts and are natural formations (unlike salt evaporation ponds, which are artificial).
A salt pan forms by evaporation of a water pool, such as a lake or pond. This happens in climates where the rate of water evaporation exceeds the rate of that is, in a desert. If the water cannot drain into the ground, it remains on the surface until it evaporates, leaving behind minerals precipitated from the salt ions dissolved in the water. Over thousands of years, the minerals (usually salts) accumulate on the surface. These minerals reflect the sun's rays and often appear as white areas.
Salt pans can be dangerous. The crust of salt can conceal a quagmire of mud that can engulf a truck. The Qattara Depression in the eastern Sahara Desert contains many such traps which served as strategic barriers during World War II.
Examples
The Bonneville Salt Flats in Utah, where many land speed records have been set, are a well-known salt pan in the arid regions of the western United States.
The Etosha pan, in the Etosha National Park in Namibia, is another prominent example of a salt pan.
The Salar de Uyuni in Bolivia is the largest salt pan in the world. As of 2024, with an estimated 23 million tons, Bolivia holds about 22% of the world's known lithium resources (105 million tons); most of those are in the Salar de Uyuni. The large area, clear skies, and exceptional flatness of the surface make the Salar an ideal object for calibrating the altimeters of Earth observation satellites.
Parts of Rann of Kutch (India) are salt marsh in the wet season and salt pan in the dry season.
| Physical sciences | Hydrology | Earth science |
390757 | https://en.wikipedia.org/wiki/Sprain | Sprain | A sprain is a soft tissue injury of the ligaments within a joint, often caused by a sudden movement abruptly forcing the joint to exceed its functional range of motion. Ligaments are tough, inelastic fibers made of collagen that connect two or more bones to form a joint and are important for joint stability and proprioception, which is the body's sense of limb position and movement. Sprains may be mild (first degree), moderate (second degree), or severe (third degree), with the latter two classes involving some degree of tearing of the ligament. Sprains can occur at any joint but most commonly occur in the ankle, knee, or wrist. An equivalent injury to a muscle or tendon is known as a strain.
The majority of sprains are mild, causing minor swelling and bruising that can be resolved with conservative treatment, typically summarized as RICE: rest, ice, compression, elevation. However, severe sprains involve complete tears, ruptures, or avulsion fractures, often leading to joint instability, severe pain, and decreased functional ability. These sprains require surgical fixation, prolonged immobilization, and physical therapy.
Signs and symptoms
Pain
Swelling
Bruising or hematoma caused by broken blood vessels within the injured ligament
Joint instability
Difficulty with bearing weight
Decreased functional ability or range of motion of the injured joint
Ligament rupture may cause a cracking or popping sound at the time of injury
Knowing the signs and symptoms of a sprain can be helpful in differentiating the injury from a strain or simple fracture. Strains typically present with pain, cramping, muscle spasm, and muscle weakness, and fractures typically present with bone tenderness, especially when bearing weight.
Causes
Acute sprains typically occur when the joint is abruptly forced beyond its functional range of motion, often in the setting of trauma or sports injuries. The most common cause of sprains in general is repetitive movements (overuse).
Mechanism
Ligaments are collagen fibers that connect bones together, providing passive stabilization to a joint. These fibers can be found in various organizational patterns (parallel, oblique, spiral, etc.) depending on the function of the joint involved. Ligaments can be extra-capsular (located outside the joint capsule), capsular (continuation of the joint capsule), or intra-articular (located within a joint capsule). The location has important implications for healing as blood flow to intra-articular ligaments is diminished compared to extra-capsular or capsular ligaments.
Collagen fibers have about a 4% elastic zone where fibers stretch out with increased load on the joint. However, exceeding this elastic limit causes a rupture of fibers, leading to a sprain. It is important to recognize that ligaments adapt to training by increasing the cross-sectional area of fibers. When a ligament is immobilized, the ligament has been shown to rapidly weaken. Normal daily activity is important for maintaining about 80–90% of the mechanical properties of a ligament.
Risk factors
Fatigue and overuse
High-intensity contact sports
Environmental factors
Poor conditioning or equipment
Age and genetic predisposition to ligament injuries
Lack of stretching or "warming up", which when performed properly increases blood flow and joint flexibility
Diagnosis
Sprains can often be diagnosed clinically based on the patient's signs and symptoms, mechanism of injury, and physical examination. However, x-rays can be obtained to help identify fractures, especially in cases of tenderness or bone pain at the injured site. In some instances, particularly if the healing process is prolonged or a more serious injury is suspected, magnetic resonance imaging (MRI) is performed to look at the surrounding soft tissue and ligaments.
Classification
First degree sprain (mild) – There is minor stretching and structural damage to the ligament, leading to mild swelling and bruising. Patients typically present without joint instability or decreased range of motion of the joint.
Second degree sprain (moderate) – There is a partial tear of the affected ligament. Patients typically experience moderate swelling, tenderness, and some instability of the joint. There may be some difficulty bearing weight on the affected joint.
Third degree sprain (severe) – There is a complete rupture or tear of the ligament, sometimes avulsing a piece of bone. Patients typically experience severe joint instability, pain, bruising, swelling, and inability to apply weight to the joint.
Joints involved
Although any joint can experience a sprain, some of the more common injuries include the following:
Ankle - Sprains most commonly occur at the ankle and can take longer to heal than ankle bone fractures. Most sprained ankles usually occur in the lateral ligaments on the outside of the ankle. Common causes include walking on uneven surfaces or during contact sports. See sprained ankle or high ankle sprain for more details.
Inversion Ankle Sprain - injury that occurs when ankle rolls inward
Eversion Ankle Sprain - injury that occurs when ankle rolls outward
Toes
Turf toe (metatarsophalangeal joint sprain) - forced hyperextension of the big toe upwards, especially during sports (initiating a sprint on a hard surface)
Knee - Sprains commonly occur at the knee, especially following intense pivoting on a planted leg during contact sports (American football, football, basketball, pole vaulting, softball, baseball and some styles of martial arts).
Anterior cruciate ligament (ACL) injury
Posterior cruciate ligament (PCL) injury
Medial collateral ligament (MCL) injury
Lateral collateral ligament (LCL) injury
Superior Tibiofibular Joint Sprain - typically caused by a twisting injury to the joint connecting the tibia (shinbone) and fibula
Patellar dislocation
Fingers and wrists - Wrist sprains commonly occur, especially during a fall on an outstretched hand.
Gamekeeper's thumb (Skier's thumb) - forceful grabbing that leads to an injury to the ulnar collateral ligament (UCL) at the metacarpophalangeal (MCP) joint of the thumb, historically found in Scottish gamekeepers
Spine
Neck sprain at the cervical vertebrae
Whiplash (Traumatic Cervical Spine Syndrome) - forced hyperextension and flexion of the neck, classically found in rear-end auto accidents
Back sprain - Back sprains are one of the most common medical complaints, often caused by poor lifting mechanics and weak core muscles.
Treatment
Treatment of sprains usually involves incorporating conservative measures to reduce the signs and symptoms of sprains, surgery to repair severe tears or ruptures, and rehabilitation to restore function to the injured joint. Although most sprains can be managed without surgery, severe injuries may require tendon grafting or ligament repair based on the individual's circumstances. The amount of rehabilitation and time needed for recovery will depend on the severity of the sprain.
A foot sprain is an injury to the ligaments that connect bones within the foot. The recovery process for a foot sprain is crucial for restoring normal function and preventing future injuries. This article outlines the general approach to foot sprain recovery, which varies depending on the severity of the injury.
Non-surgical
Depending on the mechanism of injury, joint involvement, and severity, most sprains can be treated using conservative measures following the acronym RICE within the first 24 hours of sustaining an injury. However, it is important to recognize that treatments should be individualized depending on the patient's particular injury and symptoms. Over-the-counter medications such as non-steroidal anti-inflammatory drugs (NSAIDs) can help relieve pain, and topical NSAIDs can be as effective as medications taken by mouth.
Protect: The injured site should be protected and immobilized, as there is an increased risk of recurrent injury to the affected ligaments.
Rest: The joint affected should be immobilized and bearing weight should be minimized. For example, walking should be limited in cases of sprained ankles.
Ice: Ice should be applied immediately to the sprain to reduce swelling and pain. Ice can be applied 3–4 times a day for 10–15 minutes at a time or until the swelling subsides and can be combined with a wrapping for support. Ice can also be used to numb pain but should only be applied for a short period of time (less than twenty minutes) for this purpose. Prolonged ice exposure can reduce blood flow to the injured area and slow the healing process.
Compression: Dressings, bandages, or wraps should be used to immobilize the sprain and provide support. When wrapping the injury, more pressure should be applied to the distal end of the injury and decrease in the direction of the heart. This helps circulate the blood from the extremities to the heart. Careful management of swelling through cold compression therapy is critical to the healing process by preventing further pooling of fluid in the sprained area. However, compression should not impede circulation of the limb.
Elevation: Keeping the sprained joint elevated (in relation to the rest of the body) can minimize swelling.
Other non-operative therapies including the continuous passive motion machine (moves joint without patient exertion) and cryocuff (type of cold compress that is activated similarly to a blood pressure cuff) have been effective in reducing swelling and improving range of motion. Recent studies have shown that traction is just as effective as the RICE technique in treating ankle sprains in pediatric patients.
Functional rehabilitation
The components of an effective rehabilitation program for all sprain injuries include increasing the range of motion of the affected joint and progressive muscle strengthening exercises. After implementing conservative measures to reduce swelling and pain, mobilizing the limb within 48–72 hours following injury has been shown to promote healing by stimulating growth factors in musculoskeletal tissues linked to cellular division and matrix remodeling.
Prolonged immobilization can delay the healing of a sprain, as it usually leads to muscle atrophy and weakness. Although prolonged immobilization can have a negative effect on recovery, a study in 1996 suggest that the use of bracing can improve healing by alleviating pain and stabilizing the injury to prevent further damage to the ligament or re-injury. When using a brace, it is necessary to ensure adequate blood flow to the extremity. Ultimately, the goal of functional rehabilitation is to return the patient to full daily activities while minimizing the risk of re-injury.
| Biology and health sciences | Types | Health |
390818 | https://en.wikipedia.org/wiki/Endive | Endive | Endive () is a leaf vegetable belonging to the genus Cichorium, which includes several similar bitter-leafed vegetables. Species include Cichorium endivia (also called endive), Cichorium pumilum (also called wild endive), and Cichorium intybus (also called chicory). Chicory includes types such as radicchio, puntarelle, and Belgian endive.
There is considerable confusion between Cichorium endivia and Cichorium intybus.
Cichorium endivia
There are two main varieties of cultivated C. endivia chicon:
Curly endive, or frisée (var. crispum). This type has narrow, green, curly outer leaves. It is sometimes called chicory in the United States and is called chicorée frisée in French. Further confusion results from the fact that frisée also refers to greens lightly wilted with oil.
Escarole, or broad-leaved endive (var. latifolia), has broad, pale green leaves and is less bitter than the other varieties. Varieties or names include broad-leaved Batavian endive, grumolo, scarola, and scarole. It is eaten like other greens, sauteed, chopped into soups and stews, or as part of a green salad.
Cichorium intybus
Cichorium intybus endive is popular in Europe, and is also known as leaf chicory.
Chemical constituents
Endive is rich in many vitamins and minerals, especially in folate and vitamins A and K, and is high in fiber. It also contains kaempferol.
| Biology and health sciences | Leafy vegetables | Plants |
390883 | https://en.wikipedia.org/wiki/Electric%20locomotive | Electric locomotive | An electric locomotive is a locomotive powered by electricity from overhead lines, a third rail or on-board energy storage such as a battery or a supercapacitor. Locomotives with on-board fuelled prime movers, such as diesel engines or gas turbines, are classed as diesel–electric or gas turbine–electric and not as electric locomotives, because the electric generator/motor combination serves only as a power transmission system.
Electric locomotives benefit from the high efficiency of electric motors, often above 90% (not including the inefficiency of generating the electricity). Additional efficiency can be gained from regenerative braking, which allows kinetic energy to be recovered during braking to put power back on the line. Newer electric locomotives use AC motor-inverter drive systems that provide for regenerative braking. Electric locomotives are quiet compared to diesel locomotives since there is no engine and exhaust noise and less mechanical noise. The lack of reciprocating parts means electric locomotives are easier on the track, reducing track maintenance. Power plant capacity is far greater than any individual locomotive uses, so electric locomotives can have a higher power output than diesel locomotives and they can produce even higher short-term surge power for fast acceleration. Electric locomotives are ideal for commuter rail service with frequent stops. Electric locomotives are used on freight routes with consistently high traffic volumes, or in areas with advanced rail networks. Power plants, even if they burn fossil fuels, are far cleaner than mobile sources such as locomotive engines. The power can also come from low-carbon or renewable sources, including geothermal power, hydroelectric power, biomass, solar power, nuclear power and wind turbines. Electric locomotives usually cost 20% less than diesel locomotives, their maintenance costs are 25–35% lower, and cost up to 50% less to run.
The chief disadvantage of electrification is the high cost for infrastructure: overhead lines or third rail, substations, and control systems. The impact of this varies depending on local laws and regulations. For example, public policy in the U.S. interferes with electrification: higher property taxes are imposed on privately owned rail facilities if they are electrified. The EPA regulates exhaust emissions on locomotive and marine engines, similar to regulations on car & freight truck emissions, in order to limit the amount of carbon monoxide, unburnt hydrocarbons, nitric oxides, and soot output from these mobile power sources. Because railroad infrastructure is privately owned in the U.S., railroads are unwilling to make the necessary investments for electrification. In Europe and elsewhere, railway networks are considered part of the national transport infrastructure, just like roads, highways and waterways, so are often financed by the state. Operators of the rolling stock pay fees according to rail use. This makes possible the large investments required for the technically and, in the long-term, also economically advantageous electrification.
History
Direct current
The first known electric locomotive was built in 1837 by chemist Robert Davidson of Aberdeen, and it was powered by galvanic cells (batteries). Davidson later built a larger locomotive named Galvani, exhibited at the Royal Scottish Society of Arts Exhibition in 1841. The seven-ton vehicle had two direct-drive reluctance motors, with fixed electromagnets acting on iron bars attached to a wooden cylinder on each axle, and simple commutators. It hauled a load of six tons at four miles per hour (6 kilometers per hour) for a distance of . It was tested on the Edinburgh and Glasgow Railway in September of the following year, but the limited power from batteries prevented its general use. It was destroyed by railway workers, who saw it as a threat to their job security.
The first electric passenger train was presented by Werner von Siemens at Berlin in 1879. The locomotive was driven by a 2.2 kW, series-wound motor, and the train, consisting of the locomotive and three cars, reached a speed of 13 km/h. During four months, the train carried 90,000 passengers on a 300-meter-long (984 feet) circular track. The electricity (150 V DC) was supplied through a third insulated rail between the tracks. A contact roller was used to collect the electricity.
The world's first electric tram line opened in Lichterfelde near Berlin, Germany, in 1881. It was built by Werner von Siemens (see Gross-Lichterfelde Tramway and Berlin Straßenbahn). Volk's Electric Railway opened in 1883 in Brighton. Also in 1883, Mödling and Hinterbrühl Tram opened near Vienna in Austria. It was the first in the world in regular service powered from an overhead line. Five years later, in the U.S. electric trolleys were pioneered in 1888 on the Richmond Union Passenger Railway, using equipment designed by Frank J. Sprague.
The first electrified Hungarian railway lines were opened in 1887. Budapest (See: BHÉV): Ráckeve line (1887), Szentendre line (1888), Gödöllő line (1888), Csepel line (1912).
Much of the early development of electric locomotion was driven by the increasing use of tunnels, particularly in urban areas. Smoke from steam locomotives was noxious and municipalities were increasingly inclined to prohibit their use within their limits. The first electrically worked underground line was the City and South London Railway, prompted by a clause in its enabling act prohibiting the use of steam power. It opened in 1890, using electric locomotives built by Mather and Platt. Electricity quickly became the power supply of choice for subways, abetted by Sprague's invention of multiple-unit train control in 1897. Surface and elevated rapid transit systems generally used steam until forced to convert by ordinance.
The first use of electrification on an American main line was on a four-mile stretch of the Baltimore Belt Line of the Baltimore and Ohio Railroad (B&O) in 1895 connecting the main portion of the B&O to the new line to New York through a series of tunnels around the edges of Baltimore's downtown. Parallel tracks on the Pennsylvania Railroad had shown that coal smoke from steam locomotives would be a major operating issue and a public nuisance. Three Bo+Bo units were initially used, the EL-1 Model. At the south end of the electrified section; they coupled onto the locomotive and train and pulled it through the tunnels. Railroad entrances to New York City required similar tunnels and the smoke problems were more acute there. A collision in the Park Avenue tunnel in 1902 led the New York State legislature to outlaw the use of smoke-generating locomotives south of the Harlem River after 1 July 1908. In response, electric locomotives began operation in 1904 on the New York Central Railroad. In the 1930s, the Pennsylvania Railroad, which had introduced electric locomotives because of the NYC regulation, electrified its entire territory east of Harrisburg, Pennsylvania.
The Chicago, Milwaukee, St. Paul, and Pacific Railroad (the Milwaukee Road), the last transcontinental line to be built, electrified its lines across the Rocky Mountains and to the Pacific Ocean starting in 1915. A few East Coastlines, notably the Virginian Railway and the Norfolk and Western Railway, electrified short sections of their mountain crossings. However, by this point electrification in the United States was more associated with dense urban traffic and the use of electric locomotives declined in the face of dieselization. Diesel shared some of the electric locomotive's advantages over steam and the cost of building and maintaining the power supply infrastructure, which discouraged new installations, brought on the elimination of most main-line electrification outside the Northeast. Except for a few captive systems (e.g. the Deseret Power Railroad), by 2000 electrification was confined to the Northeast Corridor and some commuter service; even there, freight service was handled by diesel. Development continued in Europe, where electrification was widespread. 1,500 V DC is still used on some lines near France and 25 kV 50 Hz is used by high-speed trains.
Alternating current
The first practical AC electric locomotive was designed by Charles Brown, then working for Oerlikon, Zürich. In 1891, Brown had demonstrated long-distance power transmission for the International Electrotechnical Exhibition, using three-phase AC, between a hydro–electric plant at Lauffen am Neckar and the expo site at Frankfurt am Main West, a distance of 280 km. Using experience he had gained while working for Jean Heilmann on steam–electric locomotive designs, Brown observed that three-phase motors had a higher power-to-weight ratio than DC motors and, because of the absence of a commutator, were simpler to manufacture and maintain. However, they were much larger than the DC motors of the time and could not be mounted in underfloor bogies: they could only be carried within locomotive bodies. In 1896, Oerlikon installed the first commercial example of the system on the Lugano Tramway. Each 30-tonne locomotive had two motors run by three-phase 750 V 40 Hz fed from double overhead lines. Three-phase motors run at a constant speed and provide regenerative braking and are thus well suited to steeply graded routes; in 1899 Brown (by then in partnership with Walter Boveri) supplied the first main-line three-phase locomotives to the 40 km Burgdorf–Thun railway (highest point 770 metres), Switzerland. The first implementation of industrial frequency single-phase AC supply for locomotives came from Oerlikon in 1901, using the designs of Hans Behn-Eschenburg and Emil Huber-Stockar; installation on the Seebach-Wettingen line of the Swiss Federal Railways was completed in 1904. The 15 kV, 50 Hz , 48 tonne locomotives used transformers and rotary converters to power DC traction motors.
In 1894, Hungarian engineer Kálmán Kandó developed a new type 3-phase asynchronous electric drive motors and generators for electric locomotives at the Fives-Lille Company. Kandó's early 1894 designs were first applied in a short three-phase AC tramway in Évian-les-Bains (France), which was constructed between 1896 and 1898.
In 1918, Kandó invented and developed the rotary phase converter, enabling electric locomotives to use three-phase motors whilst supplied via a single overhead wire, carrying the simple industrial frequency (50 Hz) single phase AC of the high voltage national networks.
Italian railways were the first in the world to introduce electric traction for the entire length of a mainline rather than just a short stretch. The 106 km Valtellina line was opened on 4 September 1902, designed by Kandó and a team from the Ganz Works. The electrical system was three-phase at 3 kV 15 Hz. The voltage was significantly higher than used earlier and it required new designs for electric motors and switching devices. The three-phase two-wire system was used on several railways in Northern Italy and became known as "the Italian system". Kandó was invited in 1905 to undertake the management of Società Italiana Westinghouse and led the development of several Italian electric locomotives. During the period of electrification of the Italian railways, tests were made as to which type of power to use: in some sections there was a 3,600 V Hz three-phase power supply, in others there was 1,500 V DC, 3 kV DC and 10 kV AC 45 Hz supply. After WW2, 3 kV DC power was chosen for the entire Italian railway system.
A later development of Kandó, working with both the Ganz works and Societa Italiana Westinghouse, was an electro-mechanical converter, allowing the use of three-phase motors from single-phase AC, eliminating the need for two overhead wires. In 1923, the first phase-converter locomotive in Hungary was constructed on the basis of Kandó's designs and serial production began soon after. The first installation, at 16 kV 50 Hz, was in 1932 on the 56 km section of the Hungarian State Railways between Budapest and Komárom. This proved successful and the electrification was extended to Hegyeshalom in 1934.
In Europe, electrification projects initially focused on mountainous regions for several reasons: coal supplies were difficult, hydroelectric power was readily available, and electric locomotives gave more traction on steeper lines. This was particularly applicable in Switzerland, where almost all lines are electrified. An important contribution to the wider adoption of AC traction came from SNCF of France after World War II. The company had assessed the industrial-frequency AC line routed through the steep Höllental Valley, Germany, which was under French administration following the war. After trials, the company decided that the performance of AC locomotives was sufficiently developed to allow all its future installations, regardless of terrain, to be of this standard, with its associated cheaper and more efficient infrastructure. The SNCF decision, ignoring as it did the of high-voltage DC already installed on French routes, was influential in the standard selected for other countries in Europe.
The 1960s saw the electrification of many European main lines. European electric locomotive technology had improved steadily from the 1920s onwards. By comparison, the Milwaukee Road class EP-2 (1918) weighed 240 t, with a power of 3,330 kW and a maximum speed of 112 km/h; in 1935, German E 18 had a power of 2,800 kW, but weighed only 108 tons and had a maximum speed of 150 km/h. On 29 March 1955, French locomotive CC 7107 reached 331 km/h. In 1960 the SJ Class Dm 3 locomotives on Swedish Railways produced a record 7,200 kW. Locomotives capable of commercial passenger service at 200 km/h appeared in Germany and France in the same period. Further improvements resulted from the introduction of electronic control systems, which permitted the use of increasingly lighter and more powerful motors that could be fitted inside the bogies (standardizing from the 1990s onwards on asynchronous three-phase motors, fed through GTO-inverters).
In the 1980s, the development of very high-speed service brought further electrification. The Japanese Shinkansen and the French TGV were the first systems for which devoted high-speed lines were built from scratch. Similar programs were undertaken in Italy, Germany and Spain; in the United States the only new mainline service was an extension of electrification over the Northeast Corridor from New Haven, Connecticut, to Boston, Massachusetts, though new electric light rail systems continued to be built.
On 2 September 2006, a standard production Siemens electric locomotive of the Eurosprinter type ES64-U4 (ÖBB Class 1216) achieved , the record for a locomotive-hauled train, on the new line between Ingolstadt and Nuremberg. This locomotive is now employed largely unmodified by ÖBB to haul their Railjet which is however limited to a top speed of 230 km/h due to economic and infrastructure concerns.
Types
An electric locomotive can be supplied with power from
Rechargeable energy storage systems, such as a battery or ultracapacitor-powered mining locomotives.
A stationary source, such as a third rail or overhead wire.
The distinguishing design features of electric locomotives are:
The type of electrical power used, AC or DC.
The method of storing (batteries, ultracapacitors) or collecting (transmission) electrical power.
The means used to couple the traction motors to the driving wheels (drivers).
Direct and alternating current
The most fundamental difference lies in the choice of AC or DC. The earliest systems used DC, as AC was not well understood and insulation material for high voltage lines was not available. DC locomotives typically run at relatively low voltage (600 to 3,000 volts); the equipment is therefore relatively massive because the currents involved are large in order to transmit sufficient power. Power must be supplied at frequent intervals as the high currents result in large transmission system losses.
As AC motors were developed, they became the predominant type, particularly on longer routes. High voltages (tens of thousands of volts) are used because this allows the use of low currents; transmission losses are proportional to the square of the current (e.g. twice the current means four times the loss). Thus, high power can be conducted over long distances on lighter and cheaper wires. Transformers in the locomotives transform this power to a low voltage and high current for the motors. A similar high voltage, low current system could not be employed with direct current locomotives because there is no easy way to do the voltage/current transformation for DC so efficiently as achieved by AC transformers.
AC traction still occasionally uses dual overhead wires instead of single-phase lines. The resulting three-phase current drives induction motors, which do not have sensitive commutators and permit easy realisation of a regenerative brake. Speed is controlled by changing the number of pole pairs in the stator circuit, with acceleration controlled by switching additional resistors in, or out, of the rotor circuit. The two-phase lines are heavy and complicated near switches, where the phases have to cross each other. The system was widely used in northern Italy until 1976 and is still in use on some Swiss rack railways. The simple feasibility of a fail-safe electric brake is an advantage of the system, while speed control and the two-phase lines are problematic.
Rectifier locomotives, which used AC power transmission and DC motors, were common, though DC commutators had problems both in starting and at low velocities. Today's advanced electric locomotives use brushless three-phase AC induction motors. These polyphase machines are powered from GTO-, IGCT- or IGBT-based inverters. The cost of electronic devices in a modern locomotive can be up to 50% of the cost of the vehicle.
Electric traction allows the use of regenerative braking, in which the motors are used as brakes and become generators that transform the motion of the train into electrical power that is then fed back into the lines. This system is particularly advantageous in mountainous operations, as descending locomotives can produce a large portion of the power required for ascending trains. Most systems have a characteristic voltage and, in the case of AC power, a system frequency. Many locomotives have been equipped to handle multiple voltages and frequencies as systems came to overlap or were upgraded. American FL9 locomotives were equipped to handle power from two different electrical systems and could also operate as diesel–electrics.
While today's systems predominantly operate on AC, many DC systems are still in use – e.g., in South Africa and the United Kingdom (750 V and 1,500 V); Netherlands, Japan, Ireland (1,500 V); Slovenia, Belgium, Italy, Poland, Russia, Spain (3,000 V) and Washington, D.C. (750 V).
Power transmission
Electrical circuits require two connections (or for three phase AC, three connections). From the beginning, the track was used for one side of the circuit. Unlike model railroads the track normally supplies only one side, the other of the circuit being provided separately.
Overhead lines
Railways generally tend to prefer overhead lines, often called "catenaries" after the support system used to hold the wire parallel to the ground. Three collection methods are possible:
Trolley pole: a long flexible pole, which engages the line with a wheel or shoe.
Bow collector: a frame that holds a long collecting rod against the wire.
Pantograph: a hinged frame that holds the collecting shoes against the wire in a fixed geometry.
Of the three, the pantograph method is best suited for high-speed operation. Some locomotives use both overhead and third rail collection (e.g. British Rail Class 92).
In Europe, the recommended geometry and shape of pantographs are defined by standard EN 50367/IEC 60486
Third rail
Mass transit systems and suburban lines often use a third rail instead of overhead wire. It allows for smaller tunnels and lower clearance under bridges, and has advantages for intensive traffic that it is a very sturdy system, not sensitive to snapping overhead wires. Some systems use four rails, especially some lines in the London Underground. One setback for third rail systems is that level crossings become more complex, usually requiring a gap section.
The original Baltimore and Ohio Railroad electrification used a sliding pickup (a contact shoe or simply the "shoe") in an overhead channel, a system quickly found to be unsatisfactory. It was replaced by a third rail, in which a pickup rides underneath or on top of a smaller rail parallel to the main track, above ground level. There are multiple pickups on both sides of the locomotive in order to accommodate the breaks in the third rail required by trackwork. This system is preferred in subways because of the close clearances it affords.
Driving the wheels
During the initial development of railroad electrical propulsion, a number of drive systems were devised to couple the output of the traction motors to the wheels. Early locomotives often used jackshaft drives. In this arrangement, the traction motor is mounted within the body of the locomotive and drives the jackshaft through a set of gears. This system was employed because the first traction motors were too large and heavy to mount directly on the axles. Due to the number of mechanical parts involved, frequent maintenance was necessary. The jackshaft drive was abandoned for all but the smallest units when smaller and lighter motors were developed,
Several other systems were devised as the electric locomotive matured. The Buchli drive was a fully spring-loaded system, in which the weight of the driving motors was completely disconnected from the driving wheels. First used in electric locomotives from the 1920s, the Buchli drive was mainly used by the French SNCF and Swiss Federal Railways. The quill drive was also developed about this time and mounted the traction motor above or to the side of the axle and coupled to the axle through a reduction gear and a hollow shaft – the quill – flexibly connected to the driving axle. The Pennsylvania Railroad GG1 locomotive used a quill drive. Again, as traction motors continued to shrink in size and weight, quill drives gradually fell out of favor in low-speed freight locomotives. In high-speed passenger locomotives used in Europe, the quill drive is still predominant.
Another drive was the "bi-polar" system, in which the motor armature was the axle itself, the frame and field assembly of the motor being attached to the truck (bogie) in a fixed position. The motor had two field poles, which allowed a limited amount of vertical movement of the armature. This system was of limited value since the power output of each motor was limited. The EP-2 bi-polar electrics used by the Milwaukee Road compensated for this problem by using a large number of powered axles.
Modern freight electric locomotives, like their Diesel–electric counterparts, almost universally use axle-hung traction motors, with one motor for each powered axle. In this arrangement, one side of the motor housing is supported by plain bearings riding on a ground and polished journal that is integral to the axle. The other side of the housing has a tongue-shaped protuberance that engages a matching slot in the truck (bogie) bolster, its purpose being to act as a torque reaction device, as well as support. Power transfer from the motor to the axle is effected by spur gearing, in which a pinion on the motor shaft engages a bull gear on the axle. Both gears are enclosed in a liquid-tight housing containing lubricating oil. The type of service in which the locomotive is used dictates the gear ratio employed. Numerically high ratios are commonly found on freight units, whereas numerically low ratios are typical of passenger engines.
Wheel arrangements
The Whyte notation system for classifying steam locomotives is not adequate for describing the variety of electric locomotive arrangements, though the Pennsylvania Railroad applied classes to its electric locomotives as if they were steam. For example, the PRR GG1 class indicates that it is arranged like two 4-6-0 class G locomotives coupled back-to-back.
UIC classification system was typically used for electric locomotives, as it could handle the complex arrangements of powered and unpowered axles and could distinguish between coupled and uncoupled drive systems.
Battery locomotive
A battery–electric locomotive (or battery locomotive) is powered by onboard batteries; a kind of battery electric vehicle.
Such locomotives are used where a diesel or conventional electric locomotive would be unsuitable. An example is maintenance trains on underground lines when the electricity supply is turned off. Another use for battery locomotives is in industrial facilities (e.g. explosives factories, oil, and gas refineries or chemical factories) where a combustion-powered locomotive (i.e., steam- or diesel-powered) could cause a safety issue due to the risks of fire, explosion or fumes in a confined space. Battery locomotives are preferred for mine railways where gas could be ignited by trolley-powered units arcing at the collection shoes, or where excessive electrical resistance could develop in the supply or return circuits, especially due to poor contact at rail joints, and allow dangerous current leakage into the ground.
The first electric locomotive built in 1837 was a battery locomotive. It was built by chemist Robert Davidson of Aberdeen in Scotland, and it was powered by galvanic cells (batteries). Another early example was at the Kennecott Copper Mine, McCarthy, Alaska, wherein 1917 the underground haulage ways were widened to enable working by two battery locomotives of . In 1928, Kennecott Copper ordered four 700-series electric locomotives with onboard batteries. These locomotives weighed and operated on 750 volts overhead trolley wire with considerable further range whilst running on batteries. The locomotives provided several decades of service using nickel–iron battery (Edison) technology. The batteries were replaced with lead-acid batteries, and the locomotives were retired shortly afterward. All four locomotives were donated to museums, but one was scrapped. The others can be seen at the Boone and Scenic Valley Railroad, Iowa, and at the Western Railway Museum in Rio Vista, California.
The Toronto Transit Commission previously operated on the Toronto subway a battery electric locomotive built by Nippon Sharyo in 1968 and retired in 2009.
London Underground regularly operates battery–electric locomotives for general maintenance work.
, battery locomotives with 7 and 14 MWh energy capacity have been ordered by rail lines and are under development.
Supercapacitor power storage
In 2020, Zhuzhou Electric Locomotive Company, manufacturers of stored electrical power systems using supercapacitors initially developed for use in trams, announced that they were extending their product line to include locomotives.
Electric locomotives around the world
Europe
Electrification is widespread in Europe, with electric multiple units commonly used for passenger trains. Due to higher density schedules, operating costs are more dominant with respect to the infrastructure costs than in the U.S. and electric locomotives have much lower operating costs than diesel. In addition, governments were motivated to electrify their railway networks due to coal shortages experienced during the First and Second World Wars.
Diesel locomotives have less power compared to electric locomotives for the same weight and dimensions. For instance, the 2,200 kW of a modern British Rail Class 66 diesel locomotive was matched in 1927 by the electric SBB-CFF-FFS Ae 4/7 (2,300 kW), which is lighter. However, for low speeds, the tractive effort is more important than power. Diesel engines can be competitive for slow freight traffic (as it is common in Canada and the U.S.) but not for passenger or mixed passenger/freight traffic as on many European railway lines, especially where heavy freight trains must be run at comparatively high speeds (80 km/h or more).
These factors led to high degrees of electrification in most European countries. In some countries, such as Switzerland, even electric shunters are common and many private sidings are served by electric locomotives. During World War II, when materials to build new electric locomotives were not available, Swiss Federal Railways installed electric heating elements in the boilers of some steam shunters, fed from the overhead supply, to deal with the shortage of imported coal.
Recent political developments in many European countries to enhance public transit have led to another boost for electric traction. In addition, gaps in the unelectrified track are closed to avoid replacing electric locomotives by diesel for these sections. The necessary modernization and electrification of these lines are possible, due to the financing of the railway infrastructure by the state.
British electric multiple units were first introduced in the 1890s, and current versions provide public transit and there are also a number of electric locomotive classes, such as: Class 76, Class 86, Class 87, Class 90, Class 91 and Class 92.
Russia and former USSR
Russia and other countries of the former Soviet Union have a mix of 3,000 V DC and 25 kV AC for historical reasons.
The special "junction stations" (around 15 over the former USSR – Vladimir, Mariinsk near Krasnoyarsk, etc.) have wiring switchable from DC to AC. Locomotive replacement is essential at these stations and is performed together with the contact wiring switching.
Most Soviet, Czech (the USSR ordered passenger electric locomotives from Škoda), Russian and Ukrainian locomotives can operate on AC or DC only. For instance, VL80 is an AC machine, with VL10 a DC version. There were some half-experimental small series such as VL82, which could switch from AC to DC and were used in small amounts around the city of Kharkiv in Ukraine, where is no junction station at many lines. Also, the latest Russian passenger locomotive EP20 and its half-experimental predecessor EP10 are a dual system.
Historically, 3,000 V DC was used for simplicity. The first experimental track was in the Georgian mountains, then the suburban zones of the largest cities were electrified for EMUs – very advantageous due to the much better dynamic of such a train compared to the steam one, which is important for suburban service with frequent stops. Then the large mountain line between Ufa and Chelyabinsk was electrified.
For some time, electric railways were only considered to be suitable for suburban or mountain lines. In around 1950, a decision was made (according to legend, by Joseph Stalin) to electrify the highly loaded plain prairie line of Omsk-Novosibirsk. After this, electrifying the major railroads at 3,000 V DC became mainstream.
25 kV AC started in the USSR in around 1960 when the industry managed to build the rectifier-based AC-wire DC-motor locomotive (all Soviet and Czech AC locomotives were such; only the post-Soviet ones switched to electronically controlled induction motors). The first major line with AC power was Mariinsk-Krasnoyarsk-Tayshet-Zima; the lines in European Russia such as Moscow-Rostov-on-Don followed.
In the 1990s, some DC lines were rebuilt as AC to allow the usage of the huge 10 MW AC locomotive of VL85. The line around Irkutsk is one of them. The DC locomotives freed by this rebuild were transferred to the St Petersburg region.
The Trans-Siberian Railway has been partly electrified since 1929, entirely since 2002. The system is 25 kV AC 50 Hz after the junction station of Mariinsk near Krasnoyarsk, 3,000 V DC before it, and train weights are up to 6,000 tonnes.
North America
Canada
Historically, Canada has used a variety of electric locomotives, primarily for moving passengers and cargo through poorly ventilated tunnels. Electric locomotives that were in use in Canada include the St. Clair Tunnel Co. Boxcab Electric, CN Boxcab Electric, and GMD GF6C. Exo in Montreal operated ALP-45DP dual-mode electro-diesel locomotives in order to allow the locomotives to traverse the poorly ventilated Mount Royal Tunnel. The locomotives run in electric mode along the entire length of the Deux-Montagnes line and along the Mascouche line between Montreal Central Station and Ahuntsic station. The locomotives run in diesel mode for the remainder of the Mascouche line and along three other non-electrified lines. However, with the conversion of the Mount Royal Tunnel into the mainline of the Réseau express métropolitain light metro system and the permanent truncation of the Mascouche line to Ahuntsic station starting in January 2020, the locomotives are run exclusively in diesel mode.
Similar to the US the flexibility of diesel locomotives and the relatively low cost of their infrastructure has led them to prevail except where legal or operational constraints dictate the use of electricity. Leading to limited electric railway infrastructure and by extension electric locomotives operating in Canada today. As of 2021, only one example exists today, GMD SW1200MG electric locomotives operated by the Iron Ore Company of Canada for a small isolated railway hauling raw ore from their Carol Lake mine to a processing plant.
In the future Toronto's GO Transit plans to operate a fleet of new electric locomotives as a part of its Regional Express Rail initiative. The feasibility of using hydrogen fuel-cell locomotives is also being studied.
United States
Electric locomotives are used for passenger trains on Amtrak's Northeast Corridor between Washington, DC, and Boston, with a branch to Harrisburg, Pennsylvania, and on some commuter rail lines. Mass transit systems and other electrified commuter lines use electric multiple units, where each car is powered. All other long-distance passenger service and, with rare exceptions, all freight is hauled by diesel–electric locomotives.
In North America, the flexibility of diesel locomotives and the relatively low cost of their infrastructure have led them to prevail except where legal or operational constraints dictate the use of electricity. An example of the latter is the use of electric locomotives by Amtrak and commuter railroads in the Northeast Corridor. Amtrak and New Jersey Transit's New York corridor use electric locomotives, currently ALP-46s, due to the prohibition on diesel operation in Penn Station and the Hudson and East River Tunnels leading to it. Some other trains to Penn Station use dual-mode locomotives that can also operate off third-rail power in the tunnels and the station.
During the steam era, some mountainous areas were electrified but these have been discontinued. The junction between electrified and non-electrified territory is the locale of engine changes; thus Northeast Corridor trains that extend south of Washington, D.C., change locomotives there. Northeast Corridor trains used to make lengthy stops in New Haven, Connecticut, as locomotives were swapped, a delay which contributed to the decision to electrify the New Haven to Boston segment of the Northeast Corridor in 2000.
Asia
China
China has over of electrified railway. With most trunk line freight and long-distance passenger trains operated using high power electric locomotives, typically in excess of of power output. Heavy freight is hauled with extremely high power multi-section locomotives, reaching up to on the "Shen 24" series of six section electric locomotives.
India
All mainline electrified routes in India use 25 kV AC overhead electrification at 50 Hz. As of March 2017, Indian Railways haul 85% of freight and passenger traffic with electric locomotives and 45,881 km of railway lines have been electrified.
Japan
Japan has come close to complete electrification largely due to the relatively short distances and mountainous terrain, which make electric service a particularly economical investment. Additionally, the mix of freight to passenger service is weighted much more toward passenger service (even in rural areas) than in many other countries, and this has helped drive government investment into the electrification of many remote lines. However, these same factors lead operators of Japanese railways to prefer EMUs over electric locomotives. The vast majority of electric passenger service in Japan is operated with EMUs, relegating electric locomotives to freight and select long-distance services.
Australia
The Victorian Railways and New South Wales Government Railways, which pioneered electric traction in Australia in the early 20th century and continue to operate 1,500 V DC electric multiple units, have withdrawn their electric locomotives.
In both states, the use of electric locomotives on principal interurban routes proved to be a qualified success. In Victoria, because only the Gippsland line was electrified, the economic advantages of electric traction were not fully realized due to the need to change locomotives for trains that ran beyond the electrified network. The Victorian Railways L class were withdrawn from service by 1987, and the Gippsland line electrification had been dismantled by 2004.
The New South Wales 86 class locomotives introduced to NSW in 1983 had a relatively short life because the cost of maintaining the infrastructure, the need to change locomotives at the extremities of the electrified network, and higher charges levied for electricity, saw diesel locomotives take over services the electrified network.
Queensland Rail implemented electrification in the 1980s and utilises the more recent 25 kV AC technology with around 1,000 km of the narrow gauge network now electrified. It operates a fleet of electric locomotives to transport coal for export, the most recent of which the 3,000 kW (4,020 HP) 3300/3400 class.
| Technology | Rail and cable transport | null |
390905 | https://en.wikipedia.org/wiki/New%20Horizons | New Horizons | {{Infobox spaceflight
| name = New Horizons
| names_list = New Frontiers 1
| image = New Horizons Transparent.png
| image_caption = New Horizons space probe
| insignia = New Horizons - Logo2 big.png
| mission_type = Pluto/Arrokoth flyby
| operator = NASA
| COSPAR_ID =
| SATCAT =
| website =
| mission_duration = Primary: 9.5 yearsElapsed:
| manufacturer = APLSwRI
| launch_mass =
| dry_mass =
| payload_mass =
| dimensions =
| power = 245 watts
| launch_date = UTC (2:00 pm EST)
| launch_rocket = Atlas V (551) AV-010
| launch_site = Cape Canaveral, SLC41
| launch_contractor = International Launch Services
| disposal_type =
| deactivated =
| last_contact =
| orbit_eccentricity = 1.41905
| orbit_inclination = 2.23014°
| orbit_epoch = January 1, 2017 (JD 2457754.5)
| interplanetary =
| instruments_list =
| programme = New Frontiers
| previous_mission =
| next_mission = Juno}}New Horizons is an interplanetary space probe launched as a part of NASA's New Frontiers program. Engineered by the Johns Hopkins University Applied Physics Laboratory (APL) and the Southwest Research Institute (SwRI), with a team led by Alan Stern, the spacecraft was launched in 2006 with the primary mission to perform a flyby study of the Pluto system in 2015, and a secondary mission to fly by and study one or more other Kuiper belt objects (KBOs) in the decade to follow, which became a mission to 486958 Arrokoth. It is the fifth space probe to achieve the escape velocity needed to leave the Solar System.
On January 19, 2006, New Horizons was launched from Cape Canaveral Air Force Station by an Atlas V rocket directly into an Earth-and-solar escape trajectory with a speed of about . It was the fastest (average speed with respect to Earth) human-made object ever launched from Earth. It is not the fastest speed recorded for a spacecraft, which, as of 2023, is that of the Parker Solar Probe. After a brief encounter with asteroid 132524 APL, New Horizons proceeded to Jupiter, making its closest approach on February 28, 2007, at a distance of . The Jupiter flyby provided a gravity assist that increased New Horizons speed; the flyby also enabled a general test of New Horizons scientific capabilities, returning data about the planet's atmosphere, moons, and magnetosphere.
Most of the post-Jupiter voyage was spent in hibernation mode to preserve onboard systems, except for brief annual checkouts. On December 6, 2014, New Horizons was brought back online for the Pluto encounter, and instrument check-out began. On January 15, 2015, the spacecraft began its approach phase to Pluto.
On July 14, 2015, at 11:49 UTC, it flew above the surface of Pluto, which at the time was 34 AU from the Sun, making it the first spacecraft to explore the dwarf planet. In August 2016, New Horizons was reported to have traveled at speeds of more than . On October 25, 2016, at 21:48 UTC, the last recorded data from the Pluto flyby was received from New Horizons. Having completed its flyby of Pluto, New Horizons then maneuvered for a flyby of Kuiper belt object 486958 Arrokoth (then nicknamed Ultima Thule), which occurred on January 1, 2019, when it was from the Sun. In August 2018, NASA cited results by Alice on New Horizons to confirm the existence of a "hydrogen wall" at the outer edges of the Solar System. This "wall" was first detected in 1992 by the two Voyager spacecraft.
New Horizons is traveling through the Kuiper belt; it is from Earth and from the Sun as of November 2024. NASA has announced it is to extend operations for New Horizons until the spacecraft exits the Kuiper belt, which is expected to occur between 2028 and 2029.
History
In August 1992, JPL scientist Robert Staehle called Pluto discoverer Clyde Tombaugh, requesting permission to visit his planet. "I told him he was welcome to it," Tombaugh later remembered, "though he's got to go one long, cold trip." The call eventually led to a series of proposed Pluto missions leading up to New Horizons.
Stamatios "Tom" Krimigis, head of the Applied Physics Laboratory's space division, one of many entrants in the New Frontiers Program competition, formed the New Horizons team with Alan Stern in December 2000. Appointed as the project's principal investigator, Stern was described by Krimigis as "the personification of the Pluto mission". New Horizons was based largely on Stern's work since Pluto 350 and involved most of the team from Pluto Kuiper Express.
The New Horizons proposal was one of five that were officially submitted to NASA. It was later selected as one of two finalists to be subject to a three-month concept study in June 2001. The other finalist, POSSE (Pluto and Outer Solar System Explorer), was a separate but similar Pluto mission concept by the University of Colorado Boulder, led by principal investigator Larry W. Esposito, and supported by the JPL, Lockheed Martin and the University of California.
However, the APL, in addition to being supported by Pluto Kuiper Express developers at the Goddard Space Flight Center and Stanford University were at an advantage; they had recently developed NEAR Shoemaker for NASA, which had successfully entered orbit around 433 Eros earlier that year, and would later land on the asteroid to scientific and engineering fanfare.
In November 2001, New Horizons was officially selected for funding as part of the New Frontiers program. However, the new NASA Administrator appointed by the Bush administration, Sean O'Keefe, was not supportive of New Horizons and effectively canceled it by not including it in NASA's budget for 2003. NASA's Associate Administrator for the Science Mission Directorate, Ed Weiler, prompted Stern to lobby for the funding of New Horizons in hopes of the mission appearing in the Planetary Science Decadal Survey, a prioritized "wish list," compiled by the United States National Research Council, that reflects the opinions of the scientific community.
After an intense campaign to gain support for New Horizons, the Planetary Science Decadal Survey of 2003–2013 was published in the summer of 2002. New Horizons topped the list of projects considered the highest priority among the scientific community in the medium-size category; ahead of missions to the Moon, and even Jupiter. Weiler stated that it was a result that "[his] administration was not going to fight". Funding for the mission was finally secured following the publication of the report. Stern's team was finally able to start building the spacecraft and its instruments, with a planned launch in January 2006 and arrival at Pluto in 2015. Alice Bowman became Mission Operations Manager (MOM).
Mission profile New Horizons is the first mission in NASA's New Frontiers mission category, larger and more expensive than the Discovery missions but smaller than the missions of the Flagship Program. The cost of the mission, including spacecraft and instrument development, launch vehicle, mission operations, data analysis, and education/public outreach, is approximately $700 million over 15 years (2001–2016). The spacecraft was built primarily by Southwest Research Institute (SwRI) and the Johns Hopkins Applied Physics Laboratory. The mission's principal investigator is Alan Stern of the Southwest Research Institute (formerly NASA Associate Administrator).
After separation from the launch vehicle, overall control was taken by Mission Operations Center (MOC) at the Applied Physics Laboratory in Howard County, Maryland. The science instruments are operated at Clyde Tombaugh Science Operations Center (T-SOC) in Boulder, Colorado. Navigation is performed at various contractor facilities, whereas the navigational positional data and related celestial reference frames are provided by the Naval Observatory Flagstaff Station through Headquarters NASA and JPL.
KinetX is the lead on the New Horizons navigation team and is responsible for planning trajectory adjustments as the spacecraft speeds toward the outer Solar System. Coincidentally the Naval Observatory Flagstaff Station was where the photographic plates were taken for the discovery of Pluto's moon Charon. The Naval Observatory itself is not far from the Lowell Observatory where Pluto was discovered.New Horizons was originally planned as a voyage to the only unexplored planet in the Solar System. When the spacecraft was launched, Pluto was still classified as a planet, later to be reclassified as a dwarf planet by the International Astronomical Union (IAU). Some members of the New Horizons team, including Alan Stern, disagree with the IAU definition and still describe Pluto as the ninth planet. Pluto's satellites Nix and Hydra also have a connection with the spacecraft: the first letters of their names (N and H) are the initials of New Horizons. The moons' discoverers chose these names for this reason, plus Nix and Hydra's relationship to the mythological Pluto.
Mementos
In addition to the science equipment, there are nine cultural artifacts traveling with the spacecraft. These include a collection of 434,738 names stored on a compact disc, a collection of images of New Horizons project personnel on another CD, a piece of Scaled Composites's SpaceShipOne, a "Not Yet Explored" USPS stamp, and two copies of the Flag of the United States.
About of Clyde Tombaugh's ashes are aboard the spacecraft, to commemorate his discovery of Pluto in 1930. A Florida state quarter coin, whose design commemorates human exploration, is included, officially as a trim weight, as is a Maryland state quarter to honor the probe's builders. One of the science packages (a dust counter) is named after Venetia Burney, who, as a child, suggested the name "Pluto" after its discovery.
Goal
The goal of the mission is to understand the formation of the Plutonian system, the Kuiper belt, and the transformation of the early Solar System. The spacecraft collected data on the atmospheres, surfaces, interiors, and environments of Pluto and its moons. It will also study other objects in the Kuiper belt. "By way of comparison, New Horizons gathered 5,000 times as much data at Pluto as Mariner did at the Red Planet."
Some of the questions the mission attempts to answer are: What is Pluto's atmosphere made of and how does it behave? What does its surface look like? Are there large geological structures? How do solar wind particles interact with Pluto's atmosphere?
Specifically, the mission's science objectives are to:
Map the surface compositions of Pluto and Charon
Characterize the geologies and morphologies of Pluto and Charon
Characterize the neutral atmosphere of Pluto and its escape rate
Search for an atmosphere around Charon
Map surface temperatures on Pluto and Charon
Search for rings and additional satellites around Pluto
Conduct similar investigations of one or more Kuiper belt objects
Design and construction
Spacecraft subsystems
The spacecraft is comparable in size and general shape to a grand piano and has been compared to a piano glued to a cocktail bar-sized satellite dish. As a point of departure, the team took inspiration from the Ulysses spacecraft, which also carried a radioisotope thermoelectric generator (RTG) and dish on a box-in-box structure through the outer Solar System. Many subsystems and components have flight heritage from APL's CONTOUR spacecraft, which in turn had heritage from APL's TIMED spacecraft.New Horizons body forms a triangle, almost thick. (The Pioneers have hexagonal bodies, whereas the Voyagers, Galileo, and Cassini–Huygens have decagonal, hollow bodies.) A 7075 aluminium alloy tube forms the main structural column, between the launch vehicle adapter ring at the "rear", and the radio dish antenna affixed to the "front" flat side. The titanium fuel tank is in this tube. The RTG attaches with a 4-sided titanium mount resembling a gray pyramid or stepstool.
Titanium provides strength and thermal isolation. The rest of the triangle is primarily sandwich panels of thin aluminum face sheet (less than ) bonded to aluminum honeycomb core. The structure is larger than strictly necessary, with empty space inside. The structure is designed to act as shielding, reducing electronics errors caused by radiation from the RTG. Also, the mass distribution required for a spinning spacecraft demands a wider triangle.
The interior structure is painted black to equalize temperature by radiative heat transfer. Overall, the spacecraft is thoroughly blanketed to retain heat. Unlike the Pioneers and Voyagers, the radio dish is also enclosed in blankets that extend to the body. The heat from the RTG adds warmth to the spacecraft while it is in the outer Solar System. While in the inner Solar System, the spacecraft must prevent overheating, hence electronic activity is limited, power is diverted to shunts with attached radiators, and louvers are opened to radiate excess heat. While the spacecraft is cruising inactively in the cold outer Solar System, the louvers are closed, and the shunt regulator reroutes power to electric heaters.
Propulsion and attitude control New Horizons has both spin-stabilized (cruise) and three-axis stabilized (science) modes controlled entirely with hydrazine monopropellant. Additional post launch delta-v of over is provided by a internal tank. Helium is used as a pressurant, with an elastomeric diaphragm assisting expulsion. The spacecraft's on-orbit mass including fuel is over on the Jupiter flyby trajectory, but would have been only for the backup direct flight option to Pluto. Significantly, had the backup option been taken, this would have meant less fuel for later Kuiper belt operations.
There are 16 thrusters on New Horizons: four and twelve plumbed into redundant branches. The larger thrusters are used primarily for trajectory corrections, and the small ones (previously used on Cassini and the Voyager spacecraft) are used primarily for attitude control and spinup/spindown maneuvers. Two star cameras are used to measure the spacecraft attitude. They are mounted on the face of the spacecraft and provide attitude information while in spin-stabilized or 3-axis mode. In between the time of star camera readings, spacecraft orientation is provided by dual redundant miniature inertial measurement units. Each unit contains three solid-state gyroscopes and three accelerometers. Two Adcole Sun sensors provide attitude determination. One detects the angle to the Sun, whereas the other measures spin rate and clocking.
Power
A cylindrical radioisotope thermoelectric generator (RTG) protrudes in the plane of the triangle from one vertex of the triangle. The RTG provided of power at launch, and was predicted to drop approximately every year, decaying to by the time of its encounter with the Plutonian system in 2015 and will decay too far to power the transmitters in the 2030s. There are no onboard batteries since RTG output is predictable, and load transients are handled by a capacitor bank and fast circuit breakers. As of January 2019, the power output of the RTG is about .
The RTG, model "GPHS-RTG", was originally a spare from the Cassini mission. The RTG contains of plutonium-238 oxide pellets. Each pellet is clad in iridium, then encased in a graphite shell. It was developed by the U.S. Department of Energy at the Materials and Fuels Complex, a part of the Idaho National Laboratory.
The original RTG design called for of plutonium, but a unit less powerful than the original design goal was produced because of delays at the United States Department of Energy, including security activities, that delayed plutonium production. The mission parameters and observation sequence had to be modified for the reduced wattage; still, not all instruments can operate simultaneously. The Department of Energy transferred the space battery program from Ohio to Argonne in 2002 because of security concerns.
The amount of radioactive plutonium in the RTG is about one-third the amount on board the Cassini–Huygens probe when it launched in 1997. The Cassini launch had been protested by multiple organizations, due to the risk of such a large amount of plutonium being released into the atmosphere in case of an accident. The United States Department of Energy estimated the chances of a launch accident that would release radiation into the atmosphere at 1 in 350, and monitored the launch because of the inclusion of an RTG on board. It was estimated that a worst-case scenario of total dispersal of on-board plutonium would spread the equivalent radiation of 80% the average annual dosage in North America from background radiation over an area with a radius of .
Flight computer
The spacecraft carries two computer systems: the Command and Data Handling system and the Guidance and Control processor. Each of the two systems is duplicated for redundancy, for a total of four computers. The processor used for its flight computers is the Mongoose-V, a 12 MHz radiation-hardened version of the MIPS R3000 CPU. Multiple redundant clocks and timing routines are implemented in hardware and software to help prevent faults and downtime. To conserve heat and mass, spacecraft and instrument electronics are housed together in IEMs (integrated electronics modules). There are two redundant IEMs. Including other functions such as instrument and radio electronics, each IEM contains 9 boards. The software of the probe runs on Nucleus RTOS operating system.
There have been two "safing" events, that sent the spacecraft into safe mode:
On March 19, 2007, the Command and Data Handling computer experienced an uncorrectable memory error and rebooted itself, causing the spacecraft to go into safe mode. The craft fully recovered within two days, with some data loss on Jupiter's magnetotail. No impact on the subsequent mission was expected.
On July 4, 2015, there was a CPU safing event triggered by an over-assignment of commanded science operations on the craft's approach to Pluto. Fortunately, the craft was able to recover within two days without major impacts on its mission. NASA scientists therefore reduced the number of scientific operations on the craft to prevent future events, which could happen during the approach with Pluto.
Telecommunications and data handling
Communication with the spacecraft is via X band. The craft had a communication rate of at Jupiter; at Pluto's distance, a rate of approximately per transmitter was expected. Besides the low data rate, Pluto's distance also causes a latency of about 4.5 hours (one-way). The NASA Deep Space Network (DSN) dishes are used to relay commands once the spacecraft is beyond Jupiter. The spacecraft uses dual modular redundancy transmitters and receivers, and either right- or left-hand circular polarization.
The downlink signal is amplified by dual redundant 12-watt traveling-wave tube amplifiers (TWTAs) mounted on the body under the dish. The receivers are low-power designs. The system can be controlled to power both TWTAs at the same time, and transmit a dual-polarized downlink signal to the DSN that nearly doubles the downlink rate. DSN tests early in the mission with this dual polarization combining technique were successful, and the capability was declared to be operational (when the spacecraft power budget permits both TWTAs to be powered).
In addition to the high-gain antenna, there are two backup low-gain antennas and a medium-gain dish. The high-gain dish has a Cassegrain reflector layout, composite construction, of diameter providing over of gain and a half-power beam width of about a degree. The prime-focus medium-gain antenna, with a aperture and 10° half-power beam width, is mounted to the forward-facing side of the high-gain antenna's secondary reflector. The forward low-gain antenna is stacked atop the feed of the medium-gain antenna. The aft low-gain antenna is mounted within the launch adapter at the rear of the spacecraft. This antenna was used only for early mission phases near Earth, just after launch and for emergencies if the spacecraft had lost attitude control.New Horizons recorded scientific instrument data to its solid-state memory buffer at each encounter, then transmitted the data to Earth. Data storage is done on two low-power solid-state recorders (one primary, one backup) holding up to s each. Because of the extreme distance from Pluto and the Kuiper belt, only one buffer load at those encounters can be saved. This is because New Horizons would require approximately 16 months after leaving the vicinity of Pluto to transmit the buffer load back to Earth. At Pluto's distance, radio signals from the space probe back to Earth took four hours and 25 minutes to traverse 4.7 billion km of space.
Part of the reason for the delay between the gathering of and transmission of data is that all of the New Horizons instrumentation is body-mounted. In order for the cameras to record data, the entire probe must turn, and the one-degree-wide beam of the high-gain antenna was not pointing toward Earth. Previous spacecraft, such as the Voyager program probes, had a rotatable instrumentation platform (a "scan platform") that could take measurements from virtually any angle without losing radio contact with Earth. New Horizons was mechanically simplified to save weight, shorten the schedule, and improve reliability during its 15-year designed lifetime.
The Voyager 2 scan platform jammed at Saturn, and the demands of long time exposures at outer planets led to a change of plans such that the entire probe was rotated to make photos at Uranus and Neptune, similar to how New Horizons rotated.
Instruments New Horizons carries seven instruments: three optical instruments, two plasma instruments, a dust sensor and a radio science receiver/radiometer. The instruments are to be used to investigate the global geology, surface composition, surface temperature, atmospheric pressure, atmospheric temperature and escape rate of Pluto and its moons. The rated power is , though not all instruments operate simultaneously. In addition, New Horizons has an Ultrastable Oscillator subsystem, which may be used to study and test the Pioneer anomaly towards the end of the spacecraft's life.
Long-Range Reconnaissance Imager (LORRI)
The Long-Range Reconnaissance Imager (LORRI) is a long-focal-length imager designed for high resolution and responsivity at visible wavelengths. The instrument is equipped with a 1024×1024 pixel by 12-bits-per-pixel monochromatic CCD imager giving a resolution of 5 μrad (~1 arcsec). The CCD is chilled far below freezing by a passive radiator on the antisolar face of the spacecraft. This temperature differential requires insulation and isolation from the rest of the structure. The aperture Ritchey–Chretien mirrors and metering structure are made of silicon carbide to boost stiffness, reduce weight and prevent warping at low temperatures. The optical elements sit in a composite light shield and mount with titanium and fiberglass for thermal isolation. Overall mass is , with the optical tube assembly (OTA) weighing about , for one of the largest silicon-carbide telescopes flown at the time (now surpassed by Herschel). For viewing on public web sites the 12-bit per pixel LORRI images are converted to 8-bit per pixel JPEG images. These public images do not contain the full dynamic range of brightness information available from the raw LORRI images files.
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Solar Wind Around Pluto (SWAP)
Solar Wind Around Pluto (SWAP) is a toroidal electrostatic analyzer and retarding potential analyzer (RPA), that makes up one of the two instruments comprising New Horizons Plasma and high-energy particle spectrometer suite (PAM), the other being PEPSSI. SWAP measures particles of up to 6.5 keV and, because of the tenuous solar wind at Pluto's distance, the instrument is designed with the largest aperture of any such instrument ever flown.
Pluto Energetic Particle Spectrometer Science Investigation (PEPSSI)
Pluto Energetic Particle Spectrometer Science Investigation (PEPSSI) is a time of flight ion and electron sensor that makes up one of the two instruments comprising New Horizons plasma and high-energy particle spectrometer suite (PAM), the other being SWAP. Unlike SWAP, which measures particles of up to 6.5 keV, PEPSSI goes up to 1 MeV. The PEPSSI sensor has been designed to measure the mass, energy and distribution of charged particles around Pluto, and is also able to differentiate between protons, electrons, and other heavy ions.
Alice Alice is an ultraviolet imaging spectrometer that is one of two photographic instruments comprising New Horizons Pluto Exploration Remote Sensing Investigation (PERSI); the other being the Ralph telescope. It resolves 1,024 wavelength bands in the far and extreme ultraviolet (from 50–), over 32 view fields. Its goal is to determine the composition of Pluto's atmosphere. This Alice instrument is derived from another Alice aboard ESA's Rosetta spacecraft. The instrument has a mass of 4.4 kg and draws 4.4 watts of power. Its primary role is to determine the relative concentrations of various elements and isotopes in Pluto's atmosphere.
In August 2018, NASA confirmed, based on results by Alice on the New Horizons spacecraft, a "hydrogen wall" at the outer edges of the Solar System that was first detected in 1992 by the two Voyager spacecraft.
Ralph telescope
The Ralph telescope, 75 mm in aperture, is one of two photographic instruments that make up New Horizons Pluto Exploration Remote Sensing Investigation (PERSI), with the other being the Alice instrument. Ralph has two separate channels: MVIC (Multispectral Visible Imaging Camera), a visible-light CCD imager with broadband and color channels; and LEISA (Linear Etalon Imaging Spectral Array), a near-infrared imaging spectrometer. LEISA is derived from a similar instrument on the Earth Observing-1 spacecraft. Ralph was named after Alice's husband on The Honeymooners, and was designed after Alice.
On June 23, 2017, NASA announced that it has renamed the LEISA instrument to the "Lisa Hardaway Infrared Mapping Spectrometer" in honor of Lisa Hardaway, the Ralph program manager at Ball Aerospace, who died in January 2017 at age 50.
Venetia Burney Student Dust Counter (VBSDC)
The Venetia Burney Student Dust Counter (VBSDC), built by students at the University of Colorado Boulder, is operating periodically to make dust measurements. It consists of a detector panel, about , mounted on the anti-solar face of the spacecraft (the ram direction), and an electronics box within the spacecraft. The detector contains fourteen polyvinylidene difluoride (PVDF) panels, twelve science and two reference, which generate voltage when impacted. Effective collecting area is . No dust counter has operated past the orbit of Uranus; models of dust in the outer Solar System, especially the Kuiper belt, are speculative. The VBSDC is always turned on measuring the masses of the interplanetary and interstellar dust particles (in the range of nano- and picograms) as they collide with the PVDF panels mounted on the New Horizons spacecraft. The measured data is expected to greatly contribute to the understanding of the dust spectra of the Solar System. The dust spectra can then be compared with those from observations of other stars, giving new clues as to where Earth-like planets can be found in the universe. The dust counter is named for Venetia Burney, who first suggested the name "Pluto" at the age of 11. A thirteen-minute short film about the VBSDC garnered an Emmy Award for student achievement in 2006.
Radio Science Experiment (REX)
The Radio Science Experiment (REX) used an ultrastable crystal oscillator (essentially a calibrated crystal in a miniature oven) and some additional electronics to conduct radio science investigations using the communications channels. These are small enough to fit on a single card. Because there are two redundant communications subsystems, there are two, identical REX circuit boards.
Journey to Pluto
Launch
On September 24, 2005, the spacecraft arrived at the Kennedy Space Center on board a C-17 Globemaster III for launch preparations. The launch of New Horizons was originally scheduled for January 11, 2006, but was initially delayed until January 17, 2006, to allow for borescope inspections of the Atlas V's kerosene tank. Further delays related to low cloud ceiling conditions downrange, and high winds and technical difficulties—unrelated to the rocket itself—prevented launch for a further two days.
The probe finally lifted off from Pad 41 at Cape Canaveral Air Force Station, Florida, directly south of Space Shuttle Launch Complex 39, at 19:00 UTC on January 19, 2006. The Centaur second stage ignited at 19:04:43 UTC and burned for 5 minutes 25 seconds. It reignited at 19:32 UTC and burned for 9 minutes 47 seconds. The ATK Star 48B third stage ignited at 19:42:37 UTC and burned for 1 minute 28 seconds. Combined, these burns successfully sent the probe on a solar-escape trajectory at . New Horizons took only nine hours to pass the Moon's orbit. Although there were backup launch opportunities in February 2006 and February 2007, only the first twenty-three days of the 2006 window permitted the Jupiter flyby. Any launch outside that period would have forced the spacecraft to fly a slower trajectory directly to Pluto, delaying its encounter by five to six years.
The probe was launched by a Lockheed Martin Atlas V 551 rocket, with a third stage added to increase the heliocentric (escape) speed. This was the first launch of the Atlas V 551 configuration, which uses five solid rocket boosters, and the first Atlas V with a third stage. Previous flights had used zero, two, or three solid boosters, but never five. The vehicle, AV-010, weighed at lift-off, and had earlier been slightly damaged when Hurricane Wilma swept across Florida on October 24, 2005. One of the solid rocket boosters was hit by a door. The booster was replaced with an identical unit, rather than inspecting and requalifying the original.
The launch was dedicated to the memory of launch conductor Daniel Sarokon, who was described by space program officials as one of the most influential people in the history of space travel.
Inner Solar System
Trajectory corrections
On January 28 and 30, 2006, mission controllers guided the probe through its first trajectory-correction maneuver (TCM), which was divided into two parts (TCM-1A and TCM-1B). The total velocity change of these two corrections was about . TCM-1 was accurate enough to permit the cancellation of TCM-2, the second of three originally scheduled corrections. On March 9, 2006, controllers performed TCM-3, the last of three scheduled course corrections. The engines burned for 76 seconds, adjusting the spacecraft's velocity by about . Further trajectory maneuvers were not needed until September 25, 2007 (seven months after the Jupiter flyby), when the engines were fired for 15 minutes and 37 seconds, changing the spacecraft's velocity by , followed by another TCM, almost three years later on June 30, 2010, that lasted 35.6 seconds, when New Horizons had already reached the halfway point (in time traveled) to Pluto.
In-flight tests and crossing of Mars orbit
During the week of February 20, 2006, controllers conducted initial in-flight tests of three onboard science instruments, the Alice ultraviolet imaging spectrometer, the PEPSSI plasma-sensor, and the LORRI long-range visible-spectrum camera. No scientific measurements or images were taken, but instrument electronics, and in the case of Alice, some electromechanical systems were shown to be functioning correctly.
On April 7, 2006, the spacecraft passed the orbit of Mars, moving at roughly away from the Sun at a solar distance of 243 million kilometers.
Asteroid 132524 APL
Because of the need to conserve fuel for possible encounters with Kuiper belt objects subsequent to the Pluto flyby, intentional encounters with objects in the asteroid belt were not planned. After launch, the New Horizons team scanned the spacecraft's trajectory to determine if any asteroids would, by chance, be close enough for observation. In May 2006 it was discovered that New Horizons would pass close to the tiny asteroid 132524 APL on June 13, 2006. Closest approach occurred at 4:05 UTC at a distance of (around one quarter of the average Earth-Moon distance). The asteroid was imaged by Ralph (use of LORRI was not possible because of proximity to the Sun), which gave the team a chance to test Ralph capabilities, and make observations of the asteroid's composition as well as light and phase curves. The asteroid was estimated to be in diameter. The spacecraft successfully tracked the rapidly moving asteroid over June 10–12, 2006.
First Pluto sighting
The first images of Pluto from New Horizons were acquired September 21–24, 2006, during a test of LORRI. They were released on November 28, 2006. The images, taken from a distance of approximately , confirmed the spacecraft's ability to track distant targets, critical for maneuvering toward Pluto and other Kuiper belt objects.
Jupiter encounter New Horizons used LORRI to take its first photographs of Jupiter on September 4, 2006, from a distance of . More detailed exploration of the system began in January 2007 with an infrared image of the moon Callisto, as well as several black-and-white images of Jupiter itself. New Horizons received a gravity assist from Jupiter, with its closest approach at 05:43:40 UTC on February 28, 2007, when it was from Jupiter. The flyby increased New Horizons speed by , accelerating the probe to a velocity of relative to the Sun and shortening its voyage to Pluto by three years.
The flyby was the center of a four-month intensive observation campaign lasting from January to June. Being an ever-changing scientific target, Jupiter has been observed intermittently since the end of the Galileo mission in September 2003. Knowledge about Jupiter benefited from the fact that New Horizons instruments were built using the latest technology, especially in the area of cameras, representing a significant improvement over Galileo cameras, which were modified versions of Voyager cameras, which, in turn, were modified Mariner cameras. The Jupiter encounter also served as a shakedown and dress rehearsal for the Pluto encounter. Because Jupiter is much closer to Earth than Pluto, the communications link can transmit multiple loadings of the memory buffer; thus the mission returned more data from the Jovian system than it was expected to transmit from Pluto.
One of the main goals during the Jupiter encounter was observing its atmospheric conditions and analyzing the structure and composition of its clouds. Heat-induced lightning strikes in the polar regions and "waves" that indicate violent storm activity were observed and measured. The Little Red Spot, spanning up to 70% of Earth's diameter, was imaged from up close for the first time. Recording from different angles and illumination conditions, New Horizons took detailed images of Jupiter's faint ring system, discovering debris left over from recent collisions within the rings or from other unexplained phenomena. The search for undiscovered moons within the rings showed no results. Travelling through Jupiter's magnetosphere, New Horizons collected valuable particle readings. "Bubbles" of plasma that are thought to be formed from material ejected by the moon Io were noticed in the magnetotail.
Jovian moons
The four largest moons of Jupiter were in poor positions for observation; the necessary path of the gravity-assist maneuver meant that New Horizons passed millions of kilometers from any of the Galilean moons. Still, its instruments were intended for small, dim targets, so they were scientifically useful on large, distant moons. Emphasis was put on Jupiter's innermost Galilean moon, Io, whose active volcanoes shoot out tons of material into Jupiter's magnetosphere, and further. Out of eleven observed eruptions, three were seen for the first time. That of Tvashtar reached an altitude of up to . The event gave scientists an unprecedented look into the structure and motion of the rising plume and its subsequent fall back to the surface. Infrared signatures of a further 36 volcanoes were noticed. Callisto's surface was analyzed with LEISA, revealing how lighting and viewing conditions affect infrared spectrum readings of its surface water ice. Minor moons such as Amalthea had their orbit solutions refined. The cameras determined their positions, acting as "reverse optical navigation".
Outer Solar System
After passing Jupiter, New Horizons spent most of its journey towards Pluto in hibernation mode. Redundant components as well as guidance and control systems were shut down to extend their life cycle, decrease operation costs and free the Deep Space Network for other missions. During hibernation mode, the onboard computer monitored the probe's systems and transmitted a signal back to Earth; a "green" code if everything was functioning as expected or a "red" code if mission control's assistance was needed. The probe was activated for about two months a year so that the instruments could be calibrated and the systems checked. The first hibernation mode cycle started on June 28, 2007, the second cycle began on December 16, 2008, the third cycle on August 27, 2009, and the fourth cycle on August 29, 2014, after a 10-week test.New Horizons crossed the orbit of Saturn on June 8, 2008, and Uranus on March 18, 2011. After astronomers announced the discovery of two new moons in the Pluto system, Kerberos and Styx, mission planners started contemplating the possibility of the probe running into unseen debris and dust left over from ancient collisions between the moons. A study based on 18 months of computer simulations, Earth-based telescope observations and occultations of the Pluto system revealed that the possibility of a catastrophic collision with debris or dust was less than 0.3% on the probe's scheduled course. If the hazard increased, New Horizons could have used one of two possible contingency plans, the so-called SHBOTs (Safe Haven by Other Trajectories). Either the probe could have continued on its present trajectory with the antenna facing the incoming particles so the more vital systems would be protected, or it could have positioned its antenna to make a course correction that would take it just from the surface of Pluto where it was expected that the atmospheric drag would have cleaned the surrounding space of possible debris.
While in hibernation mode in July 2012, New Horizons started gathering scientific data with SWAP, PEPSSI and VBSDC. Although it was originally planned to activate just the VBSDC, other instruments were powered on in order to collect valuable heliospheric data. Before activating the other two instruments, ground tests were conducted to make sure that the expanded data gathering in this phase of the mission would not limit available energy, memory and fuel in the future and that all systems were functioning during the flyby. The first set of data was transmitted in January 2013 during a three-week activation from hibernation. The command and data handling software was updated to address the problem of computer resets.
Possible Neptune trojan targets
Other possible targets were Neptune trojans. The probe's trajectory to Pluto passed near Neptune's trailing Lagrange point (""), which may host hundreds of bodies in 1:1 resonance. In late 2013, New Horizons passed within of the high-inclination L5 Neptune trojan , which was discovered shortly before by the New Horizons KBO Search task, a survey to find additional distant objects for New Horizons to fly by after its 2015 encounter with Pluto. At that range, would have been bright enough to be detectable by New Horizons LORRI instrument; however, the New Horizons team eventually decided that they would not target for observations because the preparations for the Pluto approach took precedence. On August 25, 2014, New Horizons crossed the orbit of Neptune, exactly 25 years after the planet was visited by the Voyager 2 probe. This was the last major planet orbit crossing before the Pluto flyby. At the time, the spacecraft was away from Neptune and from the Sun.
Observations of Pluto and Charon 2013–14
Images from July 1 to 3, 2013, by LORRI were the first by the probe to resolve Pluto and Charon as separate objects. On July 14, 2014, mission controllers performed a sixth trajectory-correction maneuver (TCM) since its launch to enable the craft to reach Pluto. Between July 19–24, 2014, New Horizons LORRI snapped 12 images of Charon revolving around Pluto, covering almost one full rotation at distances ranging from about . In August 2014, astronomers made high-precision measurements of Pluto's location and orbit around the Sun using the Atacama Large Millimeter/submillimeter Array (ALMA), an array of radio telescopes located in Chile, to help NASA's New Horizons spacecraft accurately home in on Pluto. On December 6, 2014, mission controllers sent a signal for the craft to "wake up" from its final Pluto-approach hibernation and begin regular operations. The craft's response that it was "awake" reached Earth on December 7, 2014, at 02:30 UTC.
Pluto approach
Distant-encounter operations at Pluto began on January 4, 2015. On this date, images of the targets with the onboard LORRI imager plus the Ralph telescope were only a few pixels in width. Investigators began taking Pluto images and background starfield images to assist mission navigators in the design of course-correcting engine maneuvers that would precisely modify the trajectory of New Horizons to aim the approach.
On February 12, 2015, NASA released new images of Pluto (taken from January 25 to 31) from the approaching probe. New Horizons was more than away from Pluto when it began taking the photos, which showed Pluto and its largest moon, Charon. The exposure time was too short to see Pluto's smaller, much fainter moons.
Investigators compiled a series of images of the moons Nix and Hydra taken from January 27 through February 8, 2015, beginning at a range of . Pluto and Charon appear as a single overexposed object at the center. The right side image has been processed to remove the background starfield. The other two, even smaller moons—Kerberos and Styx—were seen on photos taken on April 25. Starting on May 11, a hazard search was performed, looking for unknown objects that could be a danger to the spacecraft, such as rings or hitherto undiscovered moons, which could then possibly be avoided by a course change. No rings or additional moons were found.
Also in regard to the approach phase during January 2015, on August 21, 2012, the team announced that they would spend mission time attempting long-range observations of the Kuiper belt object temporarily designated VNH0004 (now designated ), when the object was at a distance of from New Horizons. The object would be too distant to resolve surface features or take spectroscopy, but it would be able to make observations that cannot be made from Earth, namely a phase curve and a search for small moons. A second object was planned to be observed in June 2015, and a third in September after the flyby; the team hoped to observe a dozen such objects through 2018. On April 15, 2015, Pluto was imaged showing a possible polar cap.
Software glitch
On July 4, 2015, New Horizons experienced a software anomaly and went into safe mode, preventing the spacecraft from performing scientific observations until engineers could resolve the problem. On July 5, NASA announced that the problem was determined to be a timing flaw in a command sequence used to prepare the spacecraft for its flyby, and the spacecraft would resume scheduled science operations on July 7. The science observations lost because of the anomaly were judged to have no impact on the mission's main objectives and minimal impact on other objectives.
The timing flaw consisted of performing two tasks simultaneously—compressing previously acquired data to release space for more data, and making a second copy of the approach command sequence—that together overloaded the spacecraft's primary computer. After the overload was detected, the spacecraft performed as designed: it switched from the primary computer to the backup computer, entered safe mode, and sent a distress call back to Earth. The distress call was received the afternoon of July 4 and alerted engineers that they needed to contact the spacecraft to get more information and resolve the issue. The resolution was that the problem happened as part of preparations for the approach, and was not expected to happen again because no similar tasks were planned for the remainder of the encounter.
Pluto system encounter
The closest approach of the New Horizons spacecraft to Pluto occurred at 11:49 UTC on July 14, 2015, at a range of from the surface and from the center of Pluto. Telemetry data confirming a successful flyby and a healthy spacecraft was received on Earth from the vicinity of the Pluto system on July 15, 2015, 00:52:37 UTC, after 22 hours of planned radio silence due to the spacecraft being pointed towards the Pluto system. Mission managers estimated a one in 10,000 chance that debris could have destroyed the probe or its communication-systems during the flyby, preventing it from sending data to Earth. The first details of the encounter were received the next day, but the download of the complete data set through the 2 kbps data downlink took just over 15 months.
Objectives
The mission's science objectives were grouped in three distinct priorities. The "primary objectives" were required. The "secondary objectives" were expected to be met but were not demanded. The "tertiary objectives" were desired. These objectives could have been skipped in favor of the above objectives. An objective to measure any magnetic field of Pluto was dropped, due to mass and the expense associated with including a magnetometer on the spacecraft. Instead, SWAP and PEPSSI could indirectly detect magnetic fields around Pluto.
Primary objectives (required)
Characterize the global geology and morphology of Pluto and Charon
Map chemical compositions of Pluto and Charon surfaces
Characterize the neutral (non-ionized) atmosphere of Pluto and its escape rate
Secondary objectives (expected)
Characterize the time variability of Pluto's surface and atmosphere
Image select Pluto and Charon areas in stereo
Map the terminators (day/night border) of Pluto and Charon with high resolution
Map the chemical compositions of select Pluto and Charon areas with high resolution
Characterize Pluto's ionosphere (upper layer of the atmosphere) and its interaction with the solar wind
Search for molecular neutral species such as molecular hydrogen, hydrocarbons, hydrogen cyanide and other nitriles in the atmosphere
Search for any Charon atmosphere
Determine bolometric Bond albedos for Pluto and Charon
Map surface temperatures of Pluto and Charon
Map any additional surfaces of outermost moons: Nix, Hydra, Kerberos, and Styx
Tertiary objectives (desired)
Characterize the energetic particle environment at Pluto and Charon
Refine bulk parameters (radii, masses) and orbits of Pluto and Charon
Search for additional moons and any rings
"The New Horizons flyby of the Pluto system was fully successful, meeting and in many cases exceeding, the Pluto objectives set out for it by NASA and the National Academy of Sciences."
Flyby details
On July 14, 2015, at 11:50 UTC, New Horizons made its closest approach to Pluto, passing within 12,500 km (7,800 mi) at a speed of 13.78 km/s (49,600 km/h; 30,800 mph), while also coming as close as 28,800 km (17,900 mi) to Charon. Starting 3.2 days prior, the spacecraft mapped Pluto and Charon with 40 km (25 mi) resolution, enabling coverage of all sides. Close-range imaging was conducted twice daily to monitor for surface changes, such as snowfall or cryovolcanism. During the flyby, LORRI captured images with up to 50 m (160 ft) resolution, MVIC created four-color global maps at 1.6 km (1 mi) resolution, and LEISA obtained near-infrared hyperspectral maps at resolutions ranging from 7 km/px (4.3 mi/px) globally to 0.6 km/px (0.37 mi/px) for selected areas.
Meanwhile, Alice characterized the atmosphere, both by emissions of atmospheric molecules (airglow), and by dimming of background stars as they pass behind Pluto (occultation). During and after closest approach, SWAP and PEPSSI sampled the high atmosphere and its effects on the solar wind. VBSDC searched for dust, inferring meteoroid collision rates and any invisible rings. REX performed active and passive radio science. The communications dish on Earth measured the disappearance and reappearance of the radio occultation signal as the probe flew by behind Pluto. The results resolved Pluto's diameter (by their timing) and atmospheric density and composition (by their weakening and strengthening pattern). (Alice can perform similar occultations, using sunlight instead of radio beacons.) Previous missions had the spacecraft transmit through the atmosphere, to Earth ("downlink"). Pluto's mass and mass distribution were evaluated by the gravitational tug on the spacecraft. As the spacecraft speeds up and slows down, the radio signal exhibited a Doppler shift. The Doppler shift was measured by comparison with the ultrastable oscillator in the communications electronics.
Reflected sunlight from Charon allowed some imaging observations of the nightside. Backlighting by the Sun gave an opportunity to highlight any rings or atmospheric hazes. REX performed radiometry of the nightside.
Satellite observations New Horizons best spatial resolution of the small satellites is at Nix, at Hydra, and approximately at Kerberos and Styx. Estimates for the dimensions of these bodies are: Nix at ; Hydra at ; Kerberos at ; and Styx at .
Initial predictions envisioned Kerberos as a relatively large and massive object whose dark surface led to it having a faint reflection. This proved to be wrong as images obtained by New Horizons on July 14 and sent back to Earth in October 2015 revealed that Kerberos was smaller in size, across with a highly reflective surface suggesting the presence of relatively clean water ice similarly to the rest of Pluto's smaller moons.
Post-Pluto events
Soon after the Pluto flyby, in July 2015, New Horizons reported that the spacecraft was healthy, its flight path was within the margins, and science data of the Pluto–Charon system had been recorded. The spacecraft's immediate task was to begin returning the 6.25 gigabytes of information collected. The free-space path loss at its distance of 4.5 light-hours () is approximately 303 dB at 7 GHz. Using the high gain antenna and transmitting at full power, the signal from EIRP is +83 dBm, and at this distance, the signal reaching Earth is −220 dBm. The received signal level (RSL) using one, un-arrayed Deep Space Network antenna with 72 dBi of forward gain equals −148 dBm. Because of the extremely low RSL, it could only transmit data at 1 to 2 kilobits per second.
By March 30, 2016, about nine months after the flyby, New Horizons reached the halfway point of transmitting this data. The transfer was completed on October 25, 2016, at 21:48 UTC, when the last piece of data—part of a Pluto–Charon observation sequence by the Ralph/LEISA imager—was received by the Johns Hopkins University Applied Physics Laboratory.
As of November 2018, at a distance of from the Sun and from 486958 Arrokoth, New Horizons was heading in the direction of the constellation Sagittarius at relative to the Sun. The brightness of the Sun from the spacecraft was magnitude −18.5.
On April 17, 2021, New Horizons reached a distance of 50 AU from the Sun, while remaining fully operational.
Mission extension
The New Horizons team requested, and received, a mission extension through 2021 to explore additional Kuiper belt objects (KBOs). Funding was secured on July 1, 2016. During this Kuiper Belt Extended Mission (KEM) the spacecraft performed a close fly-by of 486958 Arrokoth and will conduct more distant observations of an additional two dozen objects, and possibly make a fly-by of another KBO.
Kuiper belt object mission
Target background
Mission planners searched for one or more additional Kuiper belt objects (KBOs) of the order of in diameter as targets for flybys similar to the spacecraft's Plutonian encounter. However, despite the large population of KBOs, many factors limited the number of possible targets. Because the flight path was determined by the Pluto flyby, and the probe only had of hydrazine propellant remaining, the object to be visited needed to be within a cone of less than a degree's width extending from Pluto. The target also needed to be within 55 AU, because beyond 55 AU, the communications link becomes too weak, and the RTG power output decays significantly enough to hinder observations. Desirable KBOs are well over in diameter, neutral in color (to contrast with the reddish Pluto), and, if possible, have a moon that imparts a wobble.
KBO Search
In 2011, mission scientists started the New Horizons KBO Search, a dedicated survey for suitable KBOs using ground telescopes. Large ground telescopes with wide-field cameras, notably the twin 6.5-meter Magellan Telescopes in Chile, the 8.2-meter Subaru Observatory in Hawaii and the Canada–France–Hawaii TelescopePluto-bound probe faces crisis (nature.com May 20, 2014) were used to search for potential targets. By participating in a citizen-science project called Ice Hunters the public helped to scan telescopic images for possible suitable mission candidates. The ground-based search resulted in the discovery of about 143 KBOs of potential interest, but none of these were close enough to the flight path of New Horizons. Only the Hubble Space Telescope was deemed likely to find a suitable target in time for a successful KBO mission. On June 16, 2014, time on Hubble was granted for a search. Hubble has a much greater ability to find suitable KBOs than ground telescopes. The probability that a target for New Horizons would be found was estimated beforehand at about 95%.
Suitable KBOs
On October 15, 2014, it was revealed that Hubble's search had uncovered three potential targets, temporarily designated PT1 ("potential target 1"), PT2 and PT3 by the New Horizons team. PT1 was eventually chosen as the target and would be named 486958 Arrokoth.
All objects had estimated diameters in the range and were too small to be seen by ground telescopes. The targets were at distances from the Sun ranging from 43 to 44 AU, which would put the encounters in the 2018–2019 period. The initial estimated probabilities that these objects were reachable within New Horizons fuel budget were 100%, 7%, and 97%, respectively. All were members of the "cold" (low-inclination, low-eccentricity) classical Kuiper belt objects, and thus were very different from Pluto.
PT1 (given the temporary designation "1110113Y" on the HST web site), the most favorably situated object, had a magnitude of 26.8, is in diameter, and was encountered in January 2019. A course change to reach it required about 35% of New Horizons available trajectory-adjustment fuel supply. A mission to PT3 was in some ways preferable, in that it is brighter and therefore probably larger than PT1, but the greater fuel requirements to reach it would have left less for maneuvering and unforeseen events.
Once sufficient orbital information was provided, the Minor Planet Center gave provisional designations to the three target KBOs: (later 486958 Arrokoth) (PT1), (PT2), and (PT3). By the fall of 2014, a possible fourth target, , had been eliminated by follow-up observations. PT2 was out of the running before the Pluto flyby.
KBO selection
On August 28, 2015, 486958 Arrokoth (then known as and nicknamed Ultima Thule) (PT1) was chosen as the flyby target. The necessary course adjustment was performed with four engine firings between October 22 and November 4, 2015. The flyby occurred on January 1, 2019, at 00:33 UTC.
Observations of other KBOs
Aside from its flyby of 486958 Arrokoth, the extended mission for New Horizons calls for the spacecraft to conduct observations of, and look for ring systems around, between 25 and 35 different KBOs. In addition, it will continue to study the gas, dust and plasma composition of the Kuiper belt before the mission extension ends in 2021.
On November 2, 2015, New Horizons imaged KBO 15810 Arawn with the LORRI instrument from . This KBO was again imaged by the LORRI instrument on April 7–8, 2016, from a distance of . The new images allowed the science team to further refine the location of 15810 Arawn to within and to determine its rotational period of 5.47 hours.
In July 2016, the LORRI camera captured some distant images of Quaoar from ; the oblique view will complement Earth-based observations to study the object's light-scattering properties.
On December 5, 2017, when New Horizons was 40.9 AU from Earth, a calibration image of the Wishing Well cluster marked the most distant image ever taken by a spacecraft (breaking the 27-year record set by Voyager 1s famous Pale Blue Dot). Two hours later, New Horizons surpassed its own record, imaging the Kuiper belt objects and from a distance of 0.50 and 0.34 AU, respectively. These were the closest images taken of a Kuiper belt object besides Pluto and Arrokoth .
The dwarf planet Haumea was observed from afar by the New Horizons spacecraft in October 2007, January 2017, and May 2020, from distances of 49 AU, 59 AU, and 63 AU, respectively. New Horizons has observed the dwarf planets Eris (2020), Haumea (2007, 2017, 2020), Makemake (2007, 2017), and Quaoar (2016, 2017, 2019), as well as the large KBOs Ixion (2016), (2016, 2017, 2019), and (2017, 2018). It also observed Neptune's largest moon Triton (which shares similarities with Pluto and Eris) in 2019.
By December 2023, New Horizons had discovered a total of about 100 KBOs, and flown close enough to about 20 of them to capture characteristics such as shape, rotational period, possible moons, and surface composition. In addition, since 2021, Canadian researchers had been able to use machine learning software to speed up identification processes of potential KBO targets for a third flyby, cutting weeks-long efforts to hours-long.
Encounter with Arrokoth
Objectives
Science objectives of the flyby included characterizing the geology and morphology of Arrokoth and mapping the surface composition (by searching for ammonia, carbon monoxide, methane, and water ice). Searches will be conducted for orbiting moonlets, a coma, rings and the surrounding environment. Additional objectives include:
Mapping the surface geology to learn how it formed and evolved
Measuring the surface temperature
Mapping the 3-D surface topography and surface composition to learn how it is similar to and different from comets such as 67P/Churyumov–Gerasimenko and dwarf planets such as Pluto
Searching for any signs of activity, such as a cloud-like coma
Searching for and studying any satellites or rings
Measuring or constraining the mass
Targeting maneuvers
Arrokoth is the first object to be targeted for a flyby that was discovered after the spacecraft was launched. New Horizons was planned to come within of Arrokoth, three times closer than the spacecraft's earlier encounter with Pluto. Images with a resolution of up to per pixel were expected.
The new mission began on October 22, 2015, when New Horizons carried out the first in a series of four initial targeting maneuvers designed to send it towards Arrokoth. The maneuver, which started at approximately 19:50 UTC and used two of the spacecraft's small hydrazine-fueled thrusters, lasted approximately 16 minutes and changed the spacecraft's trajectory by about . The remaining three targeting maneuvers took place on October 25, October 28, and November 4, 2015.
Approach phase
The craft was brought out of its hibernation at approximately 00:33 UTC on June 5, 2018 (06:12 UTC ERT, Earth-Received Time), in order to prepare for the approach phase. After verifying its health status, the spacecraft transitioned from a spin-stabilized mode to a three-axis-stabilized mode on August 13, 2018. The official approach phase began on August 16, 2018, and continued through December 24, 2018.New Horizons made its first detection of Arrokoth on August 16, 2018, from a distance of . At that time, Arrokoth was visible at magnitude 20 against a crowded stellar background in the direction of the constellation Sagittarius.
Flyby
The Core phase began a week before the encounter and continued for two days after the encounter. The spacecraft flew by the object at a speed of and within . The majority of the science data was collected within 48 hours of the closest approach in a phase called the Inner Core. Closest approach occurred January 1, 2019, at 05:33 UTC SCET at which point the probe was from the Sun. At this distance, the one-way transit time for radio signals between Earth and New Horizons was six hours. Confirmation that the craft had succeeded in filling its digital recorders occurred when data arrived on Earth ten hours later, at 15:29 UTC.
Data download
After the encounter, preliminary, high-priority data was sent to Earth on January 1 and 2, 2019. On January 9, New Horizons returned to a spin-stabilized mode to prepare sending the remainder of its data back to Earth. This download was expected to take 20 months at a data rate of 1–2 kilobits per second.
As of July 2022, approximately 10% of the data was still left to be received.
Post-Arrokoth events
In April 2020, New Horizons was used in conjunction with telescopes on Earth to take pictures of nearby stars Proxima Centauri and Wolf 359; the images from each vantage point – over 6.4 billion km (4 billion miles) apart – were compared to produce "the first demonstration of an easily observable stellar parallax."
Images taken by the LORRI camera while New Horizons was 42 to 45 AU from the Sun were used to measure the cosmic optical background, the visible light analog of the cosmic microwave background, in seven high galactic latitude fields. At that distance New Horizons saw a sky ten times darker than the sky seen by the Hubble Space Telescope because of the absence of diffuse background sky brightness from the zodiacal light in the inner solar system. These measurements indicate that the total amount of light emitted by all galaxies at ultraviolet and visible wavelengths may be lower than previously thought.
The spacecraft reached a distance of on April 17, 2021, at 12:42 UTC, a feat performed only four times before, by Pioneer 10, Pioneer 11, Voyager 1, and Voyager 2. Voyager 1 is the farthest spacecraft from the Sun, more than away when New Horizons reached its landmark in 2021. The support team continued to use the spacecraft in 2021 to study the heliospheric environment (plasma, dust and gas) and to study other Kuiper Belt objects.
Plans
After the spacecraft passed Arrokoth, the instruments continue to have enough power to be operational until the 2030s.
Team leader Alan Stern stated there is potential for a third flyby in the 2020s at the outer edges of the Kuiper belt. This depends on a suitable Kuiper belt object being found or confirmed close enough to the spacecraft's current trajectory. Since May 2020, the New Horizons team has been using time on the Subaru Telescope to look for suitable candidates within the spacecraft's proximity. As of June 2024, no suitable targets have been found. Beginning in fiscal year 2025, New Horizons will focus on specific heliophysics data, as stated by NASA in September 2023. It will remain available for a flyby of a different target until it leaves the Kuiper belt in 2028.New Horizons may also take a picture of Earth from its distance in the Kuiper belt, but only after completing all planned KBO flybys and imaging Uranus and Neptune. This is because pointing a camera towards Earth could cause the camera to be damaged by sunlight, as none of New Horizons' cameras have an active shutter mechanism.
Speed New Horizons has been called "the fastest spacecraft ever launched" because it left Earth at . It is also the first spacecraft launched directly into a solar escape trajectory, which requires an approximate speed while near Earth of , plus additional delta-v to cover air and gravity drag, all to be provided by the launch vehicle. As of May 2, 2024, the spacecraft is from the Sun traveling at .
However, it is not the fastest spacecraft to leave the Solar System. , this record is held by Voyager 1, traveling at relative to the Sun. Voyager 1 attained greater hyperbolic excess velocity than New Horizons due to gravity assists by Jupiter and Saturn. When New Horizons reaches the distance of , it will be traveling at about , around slower than Voyager 1 at that distance. The Parker Solar Probe can also be measured as the fastest object, because of its orbital speed relative to the Sun at perihelion: . Because it remains in solar orbit, its specific orbital energy relative to the Sun is lower than New Horizons and other artificial objects escaping the Solar System.New Horizons Star 48B third stage is also on a hyperbolic escape trajectory from the Solar System and reached Jupiter before the New Horizons spacecraft; it was expected to cross Pluto's orbit on October 15, 2015. Because it was not in controlled flight, it did not receive the correct gravity assist and passed within of Pluto. The Centaur second stage did not achieve solar escape velocity and remains in a heliocentric orbit.
Gallery
Images of the launch
Videos
| Technology | Unmanned spacecraft | null |
391020 | https://en.wikipedia.org/wiki/Sirenidae | Sirenidae | Sirenidae, the sirens, are a family of neotenic aquatic salamanders. Family members have very small fore limbs and lack hind limbs altogether. In one species, the skeleton in their fore limbs is made of only cartilage. In contrast to most other salamanders, they have external gills bunched together on the neck in both larval and adult states. Sirens are found only in the Southeastern United States and northern Mexico.
Although they are primarily carnivorous, they are the only salamanders observed eating plant material.
Description
Sirens are quite distinct from other salamanders, and in some classifications they form their own suborder, Sirenoidea, or as a completely distinct order (Meantes or Trachystomata). Genetic analysis variously places them as the sister to other Salamandroidea or as sister to all other salamanders. Many of their unique characteristics seem to be partly primitive and partly derivative.
Sirens are generally eel-like in form, with two tiny, but otherwise fully developed, fore limbs. They range from in length. They are neotenic, although the larval gills are small and functionless at first, and only adults have fully developed (but inefficient) gills. They are obligate air-breathers with well developed lungs. Proving they likely evolved from a terrestrial ancestor with an aquatic larval stage. Like amphiumas, they are able to cross land on rainy nights.
These amphibians are omnivorous, feeding mainly on worms, small snails, shrimps, and filamentous algae. They are notable among salamanders (and most amphibians, aside from a few frog species) due to their semi-herbivorous habits. Their jaws possess sharp-edged keratinised and toothless ridges, like many anuran tadpoles, but the coronoid bone in the lower jaw and the vomer and palatine bone in the upper jaw have patches of monocuspid and unbladed teeth arranged in polystichous patterns.
If the conditions of a water source are unsuitable, a larva will shrink its gills to mere stumps, and these may not function at all. They are also able to burrow into mud of drying ponds and encase themselves with a cocoon of mucus to survive periods of drought. During such periods, they breathe with their small but functional lungs.
Unlike other salamanders, an interventricular septum is present in the heart. At least two of the species can produce vocalizations.
The structure of sirens' reproductive systems suggests they employ external fertilization. This has finally been confirmed in captive breeding experiments, showcasing that males also engage in parental care, building nests for their offspring. Parental care among sirens is paternal due to external fertilization. In S. intermedia males circle around females and may rub or bite her flank region. Both male and female will go on their backs and turn. It is assumed here where the female spawns and the male fertilizes her eggs. After the courtship is over, the female leaves and the male guards the eggs. Males could potentially guard more than one brood, but they are known to bite females who enter a nesting site. Paternal care has also been observed in Cryptobranchoidea, the other suborder with external fertilization. This is critical to phylogeny, as Salamandroidea, the third suborder, use internal fertilization which may be pair with maternal care, meaning that sirens are one of the oldest groups of salamanders.
The combined biomass of Siren intermedia species in a Texas pond exceeded the total biomass of the pond's seven species of fish.
Taxonomy
The siren family (Sirenidae) is subdivided into five genera, three extinct, and two extant with two and three extant species, respectively:
Genus †Habrosaurus Gilmore 1928
†H. dilatus Gilmore 1928
†H. prodilatus Gardner 2003
Genus †Kababisha Evans et al. 1996
†K. humarensis Evans et al. 1996
†K. sudanensis Evans et al. 1996
Genus †Noterpeton Rage et al. 1993
†Noterpeton bolivianum Rage et al. 1993
Genus Pseudobranchus Gray 1825 dwarf sirens
†P. robustus Goin and Auffenberg 1955
†P. vetustus
P. axanthus Netting & Goin 1942 southern dwarf siren
P. striatus LeConte 1824 northern dwarf siren
Genus Siren Österdam 1766 sirens
†S. dunni Goin and Auffenberg 1957
†S. hesterna
†S. miotexana
†S. simpsoni
S. intermedia Barnes 1826 lesser siren
S. lacertina Linnaeus, 1766 greater siren
S. nettingi Goin, 1942 western siren
S. reticulata Graham, Kline, Steen & Kelehear, 2018 reticulated siren
S. sphagnicola Fedler, Enge, & Moler, 2023 seepage siren
| Biology and health sciences | Salamanders and newts | Animals |
391037 | https://en.wikipedia.org/wiki/Northern%20mockingbird | Northern mockingbird | The northern mockingbird (Mimus polyglottos) is a mockingbird commonly found in North America, of the family Mimidae. The species is also found in some parts of the Caribbean, as well as on the Hawaiian Islands. It is typically a permanent resident across much of its range, but northern mockingbirds may move farther south during inclement weather or prior to the onset of winter. The northern mockingbird has gray to brown upper feathers and a paler belly. Its tail and wings have white patches which are visible in flight.
The species is known for its ability to mimic bird calls and other types of sound, including artificial and electronic noises. Studies have shown its ability to identify individual humans and treat them differently based on learned threat assessments. It is an omnivore and consumes fruit, invertebrates, and small vertebrates. It is often found in open areas, open woodlands and forest edges, and is quite common in urbanized areas. The species breeds from southeastern Canada throughout the United States to the Greater Antilles. It is listed as a species of least concern by the International Union for Conservation of Nature.
The mockingbird is influential in United States culture, being the state bird of five states, appearing in book titles, songs and lullabies, and making other appearances in popular culture.
Taxonomy
Swedish zoologist Carl Linnaeus first described this species in his Systema Naturae in 1758 as a species of thrush, Turdus polyglottos. Its current genus name, Mimus, is Latin for "mimic" and the specific polyglottos, is from Ancient Greek , "harmonious", from polus, "many", and glossa, "tongue", representing its outstanding ability to mimic various sounds. The northern mockingbird is considered to be conspecific with the tropical mockingbird (Mimus gilvus).
This species is categorized as the northern mockingbird as the closest living relative to M. gilvus.
Subspecies
There are three recognized subspecies for the northern mockingbird. There have been proposed races from the Bahamas and Haiti placed under the orpheus section.
M. p. polyglottos (Linnaeus, 1758): generally found in the eastern portion of North America ranging from Nova Scotia to Nebraska, to as far south as Texas and Florida.
M. p. leucopterus, the western mockingbird (Vigors, 1839): generally found in the western portion of North America ranging from northwestern Nebraska and western Texas to the Pacific coast, and south to Mexico (the Isthmus of Tehuantepec), and Socorro Island. It is larger than M. p. polyglottos and has a slightly shorter tail, upperparts are more buff and paler, underparts have a stronger buff pigment.
M. p. orpheus (Linnaeus, 1758) from the Bahamas to the Greater Antilles, also the Cayman and Virgin Islands. Similar to M. m. polyglottos except smaller, a paler shade of gray on its back, and underparts with practically little, if any buff at all.
Description
The northern mockingbird is a medium-sized mimid that has long legs and tail. Males and females look alike. Its upper parts are colored gray, while its underparts have a white or whitish-gray color.
It has parallel wing bars on the half of the wings connected near the white patch giving it a distinctive appearance in flight. The black central rectrices and typical white lateral rectrices are also noticeable in flight. The iris is usually a light green-yellow or a yellow, but there have been instances of an orange color. The bill is black with a brownish black appearance at the base. The juvenile appearance is marked by its streaks on its back, distinguished spots and streaks on its chest, and a gray or grayish-green iris.
Northern mockingbirds measure from including a tail almost as long as its body. The wingspan can range from and body mass is from . Males tend to be slightly larger than females. Among standard measurements, the wing chord is , the tail is , the culmen is and the tarsus is .
The northern mockingbird's lifespan is observed to be up to 8 years, but captive birds can live up to 20 years.
Distribution and habitat
The mockingbird's breeding range is from the Maritime provinces westwards to British Columbia, practically the entire continental United States south of the northern Plains states and Pacific Northwest, the Greater Antilles, and the majority of Mexico to eastern Oaxaca and Veracruz. The mockingbird is generally a year-round resident of its range, but the birds that live in the northern portion of its range have been noted farther south during the winter season. Sightings of the mockingbird have also been recorded in Hawaii (where it was introduced in the 1920s), southeastern Alaska, and three times as transatlantic vagrants in Britain, most recently in Exmouth, Devon, UK in February and March 2021. The mockingbird is thought to be at least partly migratory in the northern portions of its range, but the migratory behavior is not well understood.
In the 19th century, the range of the mockingbird expanded northward towards provinces such as Nova Scotia and Ontario and states such as Massachusetts, although the sightings were sporadic. Within the first five decades of the 20th century, regions that received an influx of mockingbirds were Maine, Vermont, Ohio, Iowa, and New York. In western states such as California, the population was restricted to the Lower Sonoran Desert regions but by the 1970s the mockingbird was residential in most counties. Islands that saw introductions of the mockingbird include Bermuda (in which it failed), Barbados, St. Helena, Socorro Island, the Cayman Islands and Tahiti.
The mockingbird's habitat varies by location, but it prefers open areas with sparse vegetation. In the eastern regions, suburban and urban areas such as parks and gardens are frequent residential areas. It has an affinity for mowed lawns with shrubs within proximity for shade and nesting. In western regions, desert scrub and chaparral are among its preferred habitats. When foraging for food, it prefers short grass. This bird does not nest in densely forested areas, and generally resides in the same habitats year round.
Behavior
Diet
The northern mockingbird is an omnivore. The birds' diet consists of arthropods (such as spiders, grasshoppers, wasps, bees, ants, beetles, butterflies, and caterpillars), earthworms, berries, fruits, seeds, and occasionally flowers, small crustaceans, and lizards (Anolis spp.). Mockingbirds can drink from puddles, river and lake edges, or dew and rain droplets that amass onto plants. Adult mockingbirds also have been seen drinking sap from the cuts on recently pruned trees. Its diet heavily consists of animal prey during the breeding season, but takes a drastic shift to fruits during the fall and winter. The drive for fruits amid winter has been noted for the geographic expansion of the mockingbird, and in particular, the fruit of Rosa multiflora, a favorite of the birds, is a possible link. Mockingbirds also eat garden fruits such as tomatoes, apples, and berries (like blackberries, raspberries, other bramble fruits, holly berries, mulberries, and dogwood), as well as grapes and figs.
These birds forage on the ground or in vegetation; they also fly down from a perch to capture food. While foraging, they frequently spread their wings in a peculiar two-step motion to display the white patches. There is disagreement among ornithologists over the purpose of this behavior, with hypotheses ranging from deceleration to intimidation of predators or prey.
Breeding
Both the male and female of the species reach sexual maturity after one year of life. The breeding season occurs in the spring and early summer. The males arrive before the beginning of the season to establish their territories. They may demonstrate or contest the edges of a territory using a boundary dance in which males, typically on the ground, face each other and hop side to side, sometimes fighting, until one flies away. The males use a series of courtship displays to attract the females to their sites. They run around the area either to showcase their territory to the females or to pursue the females. The males also engage in flight to showcase their wings. They sing and call as they perform all of these displays. The species can remain monogamous for many years, but incidents of polygyny and bigamy have been reported to occur during a single bird's lifetime.
Both the male and female are involved in the nest building. The male does most of the work, while the female perches on the shrub or tree where the nest is being built to watch for predators. The nest is built approximately three to ten feet above the ground. The outer part of the nest is composed of twigs, while the inner part is lined with grasses, dead leaves, moss, or artificial fibers. The eggs are a light blue or greenish color and speckled with dots. The female lays three to five eggs, and she incubates them for nearly two weeks. Once the eggs are hatched, both the male and female will feed the chicks.
The birds aggressively defend their nests and surrounding areas against other birds and animals. When a predator is persistent, mockingbirds from neighboring territories may be summoned by distinct calls to join the defense. Other birds may gather to watch as the mockingbirds drive away the intruder. In addition to harassing domestic cats and dogs that they consider a threat, mockingbirds will at times target humans. The birds are bold, and will attack much larger birds, even hawks. One incident in Tulsa, Oklahoma involving a postal carrier resulted in the distribution of a warning letter to residents.
The northern mockingbird pairs hatch about two to four broods a year. In one breeding attempt, the northern mockingbird lays an average of four eggs. They are pale blue or greenish white with red or brown blotches, and measure about . They hatch after about 11 to 14 days of incubation by the female. After about 10 to 15 days of life, the offspring become independent.
Sexual selection
Northern mockingbirds are famous for their song repertoires. Studies have shown that males sing songs at the beginning of breeding season to attract females. Unmated males sing songs in more directions and sing more bouts than mated males. In addition, unmated males perform more flight displays than mated males. The mockingbirds usually nest several times during one breeding season. Depending on the stage of breeding and the mating status, a male mockingbird will vary his song production. The unmated male keeps close track of this change. He sings in one direction when he perceives a chance to lure a female from the nest of the mated male. Unmated males are also more likely to use elevated perches to make their songs audible farther away. Though the mockingbirds are socially monogamous, mated males have been known to sing to attract additional mates.
An observational study by Logan demonstrates that the female is continuously evaluating the quality of the male and his territory. The assessment is usually triggered by the arrival of a new male in a neighboring territory at the beginning of a new breeding season. In those cases, the mated female is constantly seen flying over both the original and the new male's territory, evaluating the qualities of both territories and exchanging calls with both males. The social mate displays aggressive behaviors towards the female, while the new male shows less aggression and sings softer songs. At the same time, both the mated male and the new male will fly over other territories to attract other females as well. Separation, mate switching and extra-pair matings do occur in northern mockingbirds.
Sex allocation
Northern mockingbirds adjust the sex ratio of their offspring according to the food availability and population density. Male offspring usually require more parental investment. There is therefore a bias for bearing the costlier sex at the beginning of a breeding season when the food is abundant. Local resource competition predicts that the parents have to share the resources with offspring that remain at the natal site after maturation. In passerine birds, like the northern mockingbird, females are more likely to disperse than males. Hence, it is adaptive to produce more dispersive sex than philopatric sex when the population density is high and the competition for local resources is intense. Since northern mockingbirds are abundant in urban environments, it is possible that the pollution and contamination in cities might affect sexual hormones and therefore play a role in offspring sex ratio.
Mating
Northern mockingbirds are socially monogamous. The sexes look alike except that the male is slightly larger than the female. Mutual mate choice is exhibited in northern mockingbirds. Both males and females prefer mates that are more aggressive towards intruders, and so exhibit greater parental investment. However, males are more defensive of their nests than females. In a population where male breeding adults outnumber female breeding adults, females have more freedom in choosing their mates. In these cases, these female breeders have the option of changing mates within a breeding season if the first male does not provide a high level of parental care, which includes feeding and nest defense. High nesting success is associated with highly aggressive males attacking intruders in the territory, and so these males are preferred by females.
Parental care
Northern mockingbirds are altricial, meaning that, when hatched, they are born relatively immobile and defenseless and therefore require nourishment for a certain duration from their parents. The young have a survival bottleneck at the nestling stage because there are higher levels of nestling predation than egg predation. The levels of belligerence exhibited by parents therefore increase once eggs hatch but there is no increase during the egg stage.
A recent study shows that both food availability and temperature affect the parental incubation of the eggs in northern mockingbirds. Increasing food availability provides the females with more time to care for the nest and perform self-maintenance. Increasing temperature, however, reduces the time the females spend at the nest and there is increased energy cost to cool the eggs. The incubation behavior is a trade-off among various environmental factors.
Mockingbird nests are also often parasitized by cowbirds. The parents are found to reject parasitic eggs at an intermediate rate. A recent study has shown that foreign eggs are more likely to be rejected from a nest later in the breeding season than from earlier in a breeding season. Early nesting hosts may not have learned the pattern and coloration of their first clutch yet, so are less likely to reject foreign eggs. There is also a seasonal threshold in terms of the overlap between the breeding seasons of the northern mockingbirds and their parasites. If the breeding season of the parasites starts later, there is less likelihood of parasitism. Hence, it pays the hosts to have relatively lower sensitivity to parasitic eggs.
Ontogeny
A laboratory observation of 38 mockingbird nestlings and fledglings (thirty-five and three, respectively) recorded the behavioral development of young mockingbirds. Notable milestones, including the eyes opening, soft vocalizations, begging, and preening, began within the first six days of life. Variation in begging and more compact movements such as perching, fear crouching, and stretching appeared by the ninth day. Wing-flashing, bathing, flight, and leaving the nest happened within seventeen days (nest leaving occurred within 11 to 13 days). Improvements of flight, walking and self-feeding took place within forty days. Agonistic behavior increased during the juvenile stages, to the extent that one of two siblings living in the same area was likely killed by the other.
Song and calls
Although many species of bird imitate the vocalizations of other birds, the northern mockingbird is the species best known in North America for doing so. Among the vocalizations it imitates are songs of the Carolina wren, northern cardinal, tufted titmouse, eastern towhee, house sparrow, wood thrush, and eastern bluebird, calls of the northern flicker and great crested flycatcher, jeers and pumphandles of the blue jay, and alarms, chups, and chirrs of the American robin. It imitates not only birds, but also other animals such as cats, dogs, frogs, and crickets and sounds from artificial items such as unoiled wheels and even car alarms. As convincing as these imitations may be to humans, they often fail to fool other birds, such as the Florida scrub jay.
The northern mockingbird's mimicry is likely to serve as a form of sexual selection through which competition between males and female choice influence a bird's song repertoire size. A 2013 study attempted to determine model selection in vocal mimics, and the data suggested that mimicry in the mockingbird resulted from the bird being genetically predisposed to learning vocalizations with acoustic characteristics such as an enlarged auditory template.
Both male and female mockingbirds sing, with the latter being generally quieter and less vocal. Male commencement of singing is in late January to February and continues into the summer and the establishing of territory into the fall. Frequency in female singing is more sporadic, as it sings less often in the summer and fall, and only sings when the male is away from the territory. The mockingbird also possesses a large song repertoire that ranges from 43 to 203 song types and the size varies by region. Repertoire sizes ranged from 14 to 150 types in Texas, and two studies of mockingbirds in Florida rounded estimates to 134 and 200, approximately. It continually expands its repertoire during its life, though it pales in comparison to mimids such as the brown thrasher.
There are four recognized calls for the mockingbird: the nest relief call, hew call, chat or chatburst, and the begging call. The hew call is mainly used by both sexes for potential nest predators, conspecific chasing, and various interactions between mates. One difference between chats and chatbursts is frequency of use, as chats are year-round and chatbursts occur in the fall. Another difference is that chatbursts appear to be used in territorial defense in the fall and chats are used by either sex when disturbed. The nest relief and begging calls are only used by the males.
Predation and threats
Adult mockingbirds can fall victim to birds of prey such as the great horned owl, screech owl and sharp-shinned hawk, though their tenacious behavior makes them less likely to be captured. Scrub jays also have killed and eaten mockingbirds. Snakes rarely capture incubating females. Fledglings have been prey to domestic cats, red-tailed hawks, and crows. Eggs and nestlings are consumed by blue jays, fish crows, American crows, red-tailed hawks, swallow-tailed kites, snakes, squirrels, and cats. Blowfly larvae and Haemoproteus have been found in Florida and Arizona populations, respectively.
Winter storms limit the expansion of mockingbirds in their range. The storms have played a role in the declining of the populations in Ohio (where it has since recovered), Michigan, Minnesota and likely in Quebec. Dry seasons also affect the mockingbird populations in Arizona.
Intelligence
In a paper published in 2009, researchers found that mockingbirds were able to recall an individual human who, earlier in the study, had approached and threatened the mockingbirds' nest. Researchers had one participant stand near a mockingbird nest and touch it, while others avoided the nest. Later, the mockingbirds recognized the intruder and exhibited defensive behavior, while ignoring the other individuals. A similar paper published in 2023 had several participants exhibit varying levels of threatening behavior towards nesting mockingbirds. During a three-day training period, "high threat" participants were instructed to touch the nest daily while accompanied by a "medium threat" participant who stood three meters away. "Low threat" participants approached the nest separately and also stood three meters away for 10 minutes. When flushed from their nests by participants during a testing period, the mockingbirds retreated further from individuals who had exhibited more threatening behavior during the training period.
Adaptation to urban habitats
The northern mockingbird is a species that is found in both urban and rural habitats. There are now more northern mockingbirds living in urban habitats than non-urban environments, so they are consequently known as an urban-positive species. Biologists have long questioned how northern mockingbirds adapt to a novel environment in cities, and whether they fall into the typical ecological traps that are common for urban-dwelling birds. A comparative study between an urban dwelling population and a rural dwelling one shows that the apparent survival is higher for individuals in the urban habitats. Lower food availability and travel costs may account for the higher mortality rate in rural habitats. Urban birds are more likely to return to the nest where they had successfully bred the previous year and avoid those where breeding success was low. One explanation for this phenomenon is that urban environments are more predictable than non-urban ones, as the site fidelity among urban birds prevents them from falling into ecological traps. Mockingbirds are also able to utilize artificial lighting in order to feed nestlings in urban areas such as residential neighborhoods into the night, in contrast to those that do not nest near those areas. The adaptation of the mockingbird in urban habitats has led it to become more susceptible to lead poisoning in Baltimore and Washington, D.C. populations.
In culture
This bird features in the title and central metaphor of the 1960 novel To Kill a Mockingbird, by Harper Lee. In that novel, mockingbirds are portrayed as innocent and generous, and two of the major characters, Atticus Finch and Miss Maudie, say it is a sin to kill a mockingbird because "they don't do one thing for us but make music for us to enjoy. They don't eat up people's gardens, don't nest in corncribs, they don't do one thing but sing their hearts out for us."
The Hunger Games franchise depicts "mockingjays," mockingbirds hybridized with jabberjays, genetically engineered birds which could memorize and repeat entire human conversations. These birds appear throughout the series as a rebellious symbol.
The traditional lullaby "Hush Little Baby" has a line that goes "Papa's gonna buy you a mockingbird".
The song of the northern mockingbird inspired many American folk songs of the mid-19th century, such as "Listen to the Mocking Bird".
Thomas Jefferson had several pet mockingbirds, including a bird named "Dick".
In the fictional Neighborhood of Make-Believe on Mister Rogers' Neighborhood, one of King Friday's "pets" is a wooden northern mockingbird on a stick, which he refers to by the scientific name Mimus polyglottos.
In 1951, Patti Page, a popular vocalist, recorded "Mockin' Bird Hill", which was sold in 10" 78 RPM format. The song reached #2 on Billboards pop music chart and reflected gentle postwar values of the period.
State bird
The northern mockingbird is the state bird of Arkansas, Florida, Mississippi, Tennessee, and Texas, and previously the state bird of South Carolina.
| Biology and health sciences | Passerida | Animals |
391141 | https://en.wikipedia.org/wiki/Basophil | Basophil | Basophils are a type of white blood cell. Basophils are the least common type of granulocyte, representing about 0.5% to 1% of circulating white blood cells. They are the largest type of granulocyte. They are responsible for inflammatory reactions during immune response, as well as in the formation of acute and chronic allergic diseases, including anaphylaxis, asthma, atopic dermatitis and hay fever. They also produce compounds that coordinate immune responses, including histamine and serotonin that induce inflammation, and heparin that prevents blood clotting, although there are less than that found in mast cell granules. Mast cells were once thought to be basophils that migrated from the blood into their resident tissues (connective tissue), but they are now known to be different types of cells.
Basophils were discovered in 1879 by German physician Paul Ehrlich, who one year earlier had found a cell type present in tissues that he termed mastzellen (now mast cells). Ehrlich received the 1908 Nobel Prize in Physiology or Medicine for his discoveries.
The name comes from the fact that these leukocytes are basophilic, i.e., they are susceptible to staining by basic dyes, as shown in the picture.
Structure
Basophils contain large cytoplasmic granules which obscure the cell nucleus under the microscope when stained. However, when unstained, the nucleus is visible and it usually has two lobes. The mast cell, another granulocyte, is similar in appearance and function. Both cell types store histamine, a chemical that is secreted by the cells when stimulated. However, they arise from different branches of hematopoiesis, and mast cells usually do not circulate in the blood stream, but instead are located in connective tissue. Like all circulating granulocytes, basophils can be recruited out of the blood into a tissue when needed.
Function
Basophils appear in many specific kinds of inflammatory reactions, particularly those that cause allergic symptoms. Basophils contain anticoagulant heparin, which prevents blood from clotting too quickly. They also contain the vasodilator histamine, which promotes blood flow to tissues. They can be found in unusually high numbers at sites of ectoparasite infection (e.g., ticks).
Like eosinophils, basophils play a role in both parasitic infections and allergies. They are found in tissues where allergic reactions are occurring and probably contribute to the severity of these reactions. Basophils have protein receptors on their cell surface that bind IgE, an immunoglobulin involved in macroparasite defense and allergy. It is the bound IgE antibody that confers a selective response of these cells to environmental substances (e.g., pollen proteins or helminth antigens).
Recent studies in mice suggest that basophils may also regulate the behavior of T cells and mediate the magnitude of the secondary immune response.
CD200
Basophil function is inhibited by CD200. Herpesvirus-6, herpesvirus-7, and herpesvirus-8 produce a CD200 homolog which also inhibits basophil function. This suggests that basophils may play a role in the immune response to these viruses. The role of basophils in the immune response to these viruses is further supported by findings that the CD200 receptor is expressed more frequently in basophils than in other circulating leukocytes.
Secretions
Basophils arise and mature in bone marrow. When activated, basophils degranulate to release histamine, proteoglycans (e.g. heparin and chondroitin), and proteolytic enzymes (e.g. elastase and lysophospholipase). They also secrete lipid mediators like leukotrienes (LTD-4), and several cytokines. Histamine and proteoglycans are pre-stored in the cell's granules while the other secreted substances are newly generated. Each of these substances contributes to inflammation. Recent evidence suggests that basophils are an important source of the cytokine, interleukin-4, perhaps more important than T cells. Interleukin-4 is considered one of the critical cytokines in the development of allergies and the production of IgE antibody by the immune system. There are other substances that can activate basophils to secrete which suggests that these cells have other roles in inflammation.
The degranulation of basophils can be investigated in vitro by using flow cytometry and the so-called basophil-activation-test (BAT). Especially, in the diagnosis of allergies including of drug reactions (e.g. induced by contrast medium), the BAT is of great impact.
Basopenia (a low basophil count) is difficult to demonstrate as the normal basophil count is so low; it has been reported in association with autoimmune urticaria (a chronic itching condition). Basophilia is also uncommon but may be seen in some forms of leukemia or lymphoma.
Clinical significance
Immunophenotyping
Basophils of mice and humans have consistent immunophenotypes, including FcεRI+, CD123, CD49b(DX-5)+, CD69+, Thy-1.2+, 2B4+, CD11bdull, CD117(c-kit)−, CD24−, CD19−, CD80−, CD14−, CD23−, Ly49c−, CD122−, CD11c−, Gr-1−, NK1.1−, B220−, CD3−, γδTCR−, αβTCR−, α4 and β4-integrin negative.
Recently, Heneberg proposed that basophils may be defined as the cellular population positive for CD13, CD44, CD54, CD63, CD69, CD107a, CD123, CD164, CD193/ CCR3, CD203c, TLR-4, and FcεRI. When activated, some additional surface markers are known to be upregulated (CD13, CD107a, CD164), or surface-exposed (CD63, and the ectoenzyme CD203c).
Allergy diagnosis
Basophils are easily isolated from venous blood and present good "indicator cells" of an IgE-mediated allergic response based on the upregulation of activation markers such as CD63 and/or CD203c upon suspect allergen stimulation. Therefore, the BAT serves to confirm IgE-mediated allergy following uncertain results from classical testing based on anamnesis, skin testing or specific IgE results. More recently, BAT has also been used for the monitoring of successful allergen immunotherapy (desensitization) to differentiate short-term desensitization versus sustained unresponsiveness to the allergen.
Etymology and pronunciation
The word basophil uses combining forms of baso- + -phil, yielding "base-loving".
Additional images
| Biology and health sciences | Circulatory system | Biology |
392019 | https://en.wikipedia.org/wiki/Mycobacterium%20tuberculosis | Mycobacterium tuberculosis | Mycobacterium tuberculosis (M. tb), also known as Koch's bacillus, is a species of pathogenic bacteria in the family Mycobacteriaceae and the causative agent of tuberculosis. First discovered in 1882 by Robert Koch, M. tuberculosis has an unusual, waxy coating on its cell surface primarily due to the presence of mycolic acid. This coating makes the cells impervious to Gram staining, and as a result, M. tuberculosis can appear weakly Gram-positive. Acid-fast stains such as Ziehl–Neelsen, or fluorescent stains such as auramine are used instead to identify M. tuberculosis with a microscope. The physiology of M. tuberculosis is highly aerobic and requires high levels of oxygen. Primarily a pathogen of the mammalian respiratory system, it infects the lungs. The most frequently used diagnostic methods for tuberculosis are the tuberculin skin test, acid-fast stain, culture, and polymerase chain reaction.
The M. tuberculosis genome was sequenced in 1998.
Microbiology
M. tuberculosis requires oxygen to grow, and is nonmotile. It divides every 18–24 hours. This is extremely slow compared with other bacteria, which tend to have division times measured in minutes (Escherichia coli can divide roughly every 20 minutes). It is a small bacillus that can withstand weak disinfectants and can survive in a dry state for weeks. Its unusual cell wall, rich in lipids such as mycolic acid and cord factor glycolipid, is likely responsible for its resistance to desiccation and is a key virulence factor.
Microscopy
Other bacteria are commonly identified with a microscope by staining them with Gram stain. However, the mycolic acid in the cell wall of M. tuberculosis does not absorb the stain. Instead, acid-fast stains such as Ziehl–Neelsen stain, or fluorescent stains such as auramine are used. Cells are curved rod-shaped and are often seen wrapped together, due to the presence of fatty acids in the cell wall that stick together. This appearance is referred to as cording, like strands of cord that make up a rope. M. tuberculosis is characterized in tissue by caseating granulomas containing Langhans giant cells, which have a "horseshoe" pattern of nuclei.
Culture
M. tuberculosis can be grown in the laboratory. Compared to other commonly studied bacteria, M. tuberculosis has a remarkably slow growth rate, doubling roughly once per day. Commonly used media include liquids such as Middlebrook 7H9 or 7H12, egg-based solid media such as Lowenstein-Jensen, and solid agar-based such as Middlebrook 7H11 or 7H10. Visible colonies require several weeks to grow on agar plates. Mycobacteria growth indicator tubes can contain a gel that emits fluorescent light if mycobacteria are grown. It is distinguished from other mycobacteria by its production of catalase and niacin. Other tests to confirm its identity include gene probes and MALDI-TOF.
Morphology
Analysis of Mycobacterium tuberculosis via scanning electron microscope shows the bacteria are in length with an average diameter of . The outer membrane and plasma membrane surface areas were measured to be and , respectively. The cell, outer membrane, periplasm, plasma membrane, and cytoplasm volumes were (= μm3), , , , and , respectively. The average total ribosome number was with ribosome density about .
Related Mycobacterium species
M. tuberculosis is part of a genetically related group of Mycobacterium species that has at least nine members:
M. tuberculosis sensu stricto
M. africanum
M. canettii
M. bovis
M. caprae
M. microti
M. pinnipedii
M. mungi
M. orygis
Pathophysiology
Humans are the only known reservoirs of M. tuberculosis. A misconception is that M. tuberculosis can be spread by shaking hands, making contact with toilet seats, sharing food or drink, or sharing toothbrushes. However, major spread is through air droplets originating from a person who has the disease either coughing, sneezing, speaking, or singing.
When in the lungs, M. tuberculosis is phagocytosed by alveolar macrophages, but they are unable to kill and digest the bacterium. Its cell wall is made of cord factor glycolipids that inhibit the fusion of the phagosome with the lysosome, which contains a host of antibacterial factors.
Specifically, M. tuberculosis blocks the bridging molecule, early endosomal autoantigen 1 (EEA1); however, this blockade does not prevent fusion of vesicles filled with nutrients. In addition, production of the diterpene isotuberculosinol prevents maturation of the phagosome. The bacteria also evades macrophage-killing by neutralizing reactive nitrogen intermediates. More recently, M. tuberculosis has been shown to secrete and cover itself in 1-tuberculosinyladenosine (1-TbAd), a special nucleoside that acts as an antacid, allowing it to neutralize pH and induce swelling in lysosomes.
In M. tuberculosis infections, PPM1A levels were found to be upregulated, and this, in turn, would impact the normal apoptotic response of macrophages to clear pathogens, as PPM1A is involved in the intrinsic and extrinsic apoptotic pathways. Hence, when PPM1A levels were increased, the expression of it inhibits the two apoptotic pathways. With kinome analysis, the JNK/AP-1 signalling pathway was found to be a downstream effector that PPM1A has a part to play in, and the apoptotic pathway in macrophages are controlled in this manner. As a result of having apoptosis being suppressed, it provides M. tuberculosis with a safe replicative niche, and so the bacteria are able to maintain a latent state for a prolonged time.
Granulomas, organized aggregates of immune cells, are a hallmark feature of tuberculosis infection. Granulomas play dual roles during infection: they regulate the immune response and minimize tissue damage, but also can aid in the expansion of infection.
The ability to construct M. tuberculosis mutants and test individual gene products for specific functions has significantly advanced the understanding of its pathogenesis and virulence factors. Many secreted and exported proteins are known to be important in pathogenesis. For example, one such virulence factor is cord factor (trehalose dimycolate), which serves to increase survival within its host. Resistant strains of M. tuberculosis have developed resistance to more than one TB drug, due to mutations in their genes. In addition, pre-existing first-line TB drugs such as rifampicin and streptomycin have decreased efficiency in clearing intracellular M. tuberculosis due to their inability to effectively penetrate the macrophage niche.
JNK plays a key role in the control of apoptotic pathways—intrinsic and extrinsic. In addition, it is also found to be a substrate of PPM1A activity, hence the phosphorylation of JNK would cause apoptosis to occur. Since PPM1A levels are elevated during M. tuberculosis infections, by inhibiting the PPM1A signalling pathways, it could potentially be a therapeutic method to kill M. tuberculosis-infected macrophages by restoring its normal apoptotic function in defence of pathogens. By targeting the PPM1A-JNK signalling axis pathway, then, it could eliminate M. tuberculosis-infected macrophages.
The ability to restore macrophage apoptosis to M. tuberculosis-infected ones could improve the current tuberculosis chemotherapy treatment, as TB drugs can gain better access to the bacteria in the niche. thus decreasing the treatment times for M. tuberculosis infections.
Symptoms of M. tuberculosis include coughing that lasts for more than three weeks, hemoptysis, chest pain when breathing or coughing, weight loss, fatigue, fever, night sweats, chills, and loss of appetite. M. tuberculosis also has the potential of spreading to other parts of the body. This can cause blood in urine if the kidneys are affected, and back pain if the spine is affected.
Strain variation
Typing of strains is useful in the investigation of tuberculosis outbreaks, because it gives the investigator evidence for or against transmission from person to person. Consider the situation where person A has tuberculosis and believes he acquired it from person B. If the bacteria isolated from each person belong to different types, then transmission from B to A is definitively disproven; however, if the bacteria are the same strain, then this supports (but does not definitively prove) the hypothesis that B infected A.
Until the early 2000s, M. tuberculosis strains were typed by pulsed field gel electrophoresis. This has now been superseded by variable numbers of tandem repeats (VNTR), which is technically easier to perform and allows better discrimination between strains. This method makes use of the presence of repeated DNA sequences within the M. tuberculosis genome.
Three generations of VNTR typing for M. tuberculosis are noted. The first scheme, called exact tandem repeat, used only five loci, but the resolution afforded by these five loci was not as good as PFGE. The second scheme, called mycobacterial interspersed repetitive unit, had discrimination as good as PFGE. The third generation (mycobacterial interspersed repetitive unit – 2) added a further nine loci to bring the total to 24. This provides a degree of resolution greater than PFGE and is currently the standard for typing M. tuberculosis. However, with regard to archaeological remains, additional evidence may be required because of possible contamination from related soil bacteria.
Antibiotic resistance in M. tuberculosis typically occurs due to either the accumulation of mutations in the genes targeted by the antibiotic or a change in titration of the drug. M. tuberculosis is considered to be multidrug-resistant (MDR TB) if it has developed drug resistance to both rifampicin and isoniazid, which are the most important antibiotics used in treatment. Additionally, extensively drug-resistant M. tuberculosis (XDR TB) is characterized by resistance to both isoniazid and rifampin, plus any fluoroquinolone and at least one of three injectable second-line drugs (i.e., amikacin, kanamycin, or capreomycin).
Genome
The genome of the H37Rv strain was published in 1998. Its size is 4 million base pairs, with 3,959 genes; 40% of these genes have had their function characterized, with possible function postulated for another 44%. Within the genome are also six pseudogenes.
Fatty acid metabolism. The genome contains 250 genes involved in fatty acid metabolism, with 39 of these involved in the polyketide metabolism generating the waxy coat. Such large numbers of conserved genes show the evolutionary importance of the waxy coat to pathogen survival. Furthermore, experimental studies have since validated the importance of a lipid metabolism for M. tuberculosis, consisting entirely of host-derived lipids such as fats and cholesterol. Bacteria isolated from the lungs of infected mice were shown to preferentially use fatty acids over carbohydrate substrates. M. tuberculosis can also grow on the lipid cholesterol as a sole source of carbon, and genes involved in the cholesterol use pathway(s) have been validated as important during various stages of the infection lifecycle of M. tuberculosis, especially during the chronic phase of infection when other nutrients are likely not available.
PE/PPE gene families. About 10% of the coding capacity is taken up by the PE/PPE gene families that encode acidic, glycine-rich proteins. These proteins have a conserved N-terminal motif, deletion of which impairs growth in macrophages and granulomas.
Noncoding RNAs. Nine noncoding sRNAs have been characterised in M. tuberculosis, with a further 56 predicted in a bioinformatics screen.
Antibiotic resistance genes. In 2013, a study on the genome of several sensitive, ultraresistant, and multiresistant M. tuberculosis strains was made to study antibiotic resistance mechanisms. Results reveal new relationships and drug resistance genes not previously associated and suggest some genes and intergenic regions associated with drug resistance may be involved in the resistance to more than one drug. Noteworthy is the role of the intergenic regions in the development of this resistance, and most of the genes proposed in this study to be responsible for drug resistance have an essential role in the development of M. tuberculosis.
Epigenome. Single-molecule real-time sequencing and subsequent bioinformatic analysis has identified three DNA methyltransferases in M. tuberculosis, Mycobacterial Adenine Methyltransferases A (MamA), B (MamB), and C (MamC). All three are adenine methyltransferases, and each are functional in some clinical strains of M. tuberculosisand not in others. Unlike DNA methyltransferases in most bacteria, which invariably methylate the adenines at their targeted sequence, some strains of M. tuberculosis carry mutations in MamA that cause partial methylation of targeted adenine bases. This occurs as intracellular stochastic methylation, where a some targeted adenine bases on a given DNA molecule are methylated while others remain unmethylated. MamA mutations causing intercellular mosaic methylation are most common in the globally successful Beijing sublineage of M. tuberculosis. Due to the influence of methylation on gene expression at some locations in the genome, it has been hypothesized that IMM may give rise to phenotypic diversity, and partially responsible for the global success of Beijing sublineage.
Evolution
The Mycobacterium tuberculosis complex (MTBC) evolved in Africa and most probably in the Horn of Africa. In addition to M. tuberculosis, the MTBC has a number of members infecting various animal species, including M. africanum, M. bovis (Dassie's bacillus), M. caprae, M. microti, M. mungi, M. orygis, and M. pinnipedii. This group may also include the M. canettii clade. These animal strains of MTBC do not strictly deserve species status, as they are all closely related and embedded in the M. tuberculosis phylogeny, but for historic reasons, they currently hold species status.
The M. canettii clade – which includes M. prototuberculosis – is a group of smooth-colony Mycobacterium species. Unlike the established members of the M. tuberculosis group, they undergo recombination with other species. The majority of the known strains of this group have been isolated from the Horn of Africa. The ancestor of M. tuberculosis appears to be M. canettii, first described in 1969.
The established members of the M. tuberculosis complex are all clonal in their spread. The main human-infecting species have been classified into seven lineages. Translating these lineages into the terminology used for spoligotyping, a very crude genotyping methodology, lineage 1 contains the East African-Indian (EAI), the Manila family of strains and some Manu (Indian) strains; lineage 2 is the Beijing group; lineage 3 includes the Central Asian (CAS) strains; lineage 4 includes the Ghana and Haarlem (H/T), Latin America-Mediterranean (LAM) and X strains; types 5 and 6 correspond to M. africanum and are observed predominantly and at high frequencies in West Africa. A seventh type has been isolated from the Horn of Africa. The other species of this complex belong to a number of spoligotypes and do not normally infect humans.
Lineages 2, 3 and 4 all share a unique deletion event (tbD1) and thus form a monophyletic group. Types 5 and 6 are closely related to the animal strains of MTBC, which do not normally infect humans. Lineage 3 has been divided into two clades: CAS-Kili (found in Tanzania) and CAS-Delhi (found in India and Saudi Arabia).
Lineage 4 is also known as the Euro-American lineage. Subtypes within this type include Latin American Mediterranean, Uganda I, Uganda II, Haarlem, X, and Congo.
A much cited study reported that M. tuberculosis has co-evolved with human populations, and that the most recent common ancestor of the M. tuberculosis complex evolved between 40,000 and 70,000 years ago. However, a later study that included genome sequences from M. tuberculosis complex members extracted from three 1,000-year-old Peruvian mummies, came to quite different conclusions. If the most recent common ancestor of the M. tuberculosis complex were 40,000 to 70,000 years old, this would necessitate an evolutionary rate much lower than any estimates produced by genomic analyses of heterochronous samples, suggesting a far more recent common ancestor of the M. tuberculosis complex as little as 6000 years ago.
An analysis of over 3000 strains of M. bovis from 35 countries suggested an Africa origin for this species.
Co-evolution with modern humans
There are currently two narratives existing in parallel regarding the age of MTBC and how it has spread and co-evolved with humans through time. One study compared the M. tuberculosis phylogeny to a human mitochondrial genome phylogeny and interpreted these as being highly similar. Based on this, the study suggested that M. tuberculosis, like humans, evolved in Africa and subsequently spread with anatomically modern humans out of Africa across the world. By calibrating the mutation rate of M. tuberculosis to match this narrative, the study suggested that MTBC evolved 40,000–70,000 years ago. Applying this time scale, the study found that the M. tuberculosis effective population size expanded during the Neolithic Demographic Transition (around 10,000 years ago) and suggested that M. tuberculosis was able to adapt to changing human populations and that the historical success of this pathogen was driven at least in part by dramatic increases in human host population density. It has also been demonstrated that after emigrating from one continent to another, a human host's region of origin is predictive of which TB lineage they carry, which could reflect either a stable association between host populations and specific M. tuberculosis lineages and/or social interactions that are shaped by shared cultural and geographic histories.
Regarding the congruence between human and M. tuberculosis phylogenies, a study relying on M. tuberculosis and human Y chromosome DNA sequences to formally assess the correlation between them, concluded that they are not congruent. Also, a more recent study which included genome sequences from M. tuberculosis complex members extracted from three 1,000-year-old Peruvian mummies, estimated that the most recent common ancestor of the M. tuberculosis complex lived only 4,000 – 6,000 years ago. The M. tuberculosis evolutionary rate estimated by the Bos et al. study is also supported by a study on Lineage 4 relying on genomic aDNA sequences from Hungarian mummies more than 200 years old. In total, the evidence thus favors this more recent estimate of the age of the MTBC most recent common ancestor, and thus that the global evolution and dispersal of M. tuberculosis has occurred over the last 4,000–6,000 years.
Among the seven recognized lineages of M. tuberculosis, only two are truly global in their distribution: Lineages 2 and 4. Among these, Lineage 4 is the most well dispersed, and almost totally dominates in the Americas. Lineage 4 was shown to have evolved in or in the vicinity of Europe, and to have spread globally with Europeans starting around the 13th century. This study also found that Lineage 4 tuberculosis spread to the Americas shortly after the European discovery of the continent in 1492, and suggests that this represented the first introduction of human TB on the continent (although animal strains have been found in human remains predating Columbus. Similarly, Lineage 4 was found to have spread from Europe to Africa during the Age of Discovery, starting in the early 15th century.
It has been suggested that ancestral mycobacteria may have infected early hominids in East Africa as early as three million years ago.
DNA fragments from M. tuberculosis and tuberculosis disease indications were present in human bodies dating from 7000 BC found at Atlit-Yam in the Levant.
Antibiotic resistance (ABR)
M. tuberculosis is a clonal organism and does not exchange DNA via horizontal gene transfer. Despite an additionally slow evolution rate, the emergence and spread of antibiotic resistance in M. tuberculosis poses an increasing threat to global public health. In 2019, the WHO reported the estimated incidence of antibiotic resistant TB to be 3.4% in new cases, and 18% in previously treated cases. Geographical discrepancies exist in the incidence rates of drug-resistant TB. Countries facing the highest rates of ABR TB China, India, Russia, and South Africa. Recent trends reveal an increase in drug-resistant cases in a number of regions, with Papua New Guinea, Singapore, and Australia undergoing significant increases.
Multidrug-resistant Tuberculosis (MDR-TB) is characterised by resistance to at least the two front-line drugs isoniazid and rifampin. MDR is associated with a relatively poor treatment success rate of 52%. Isoniazid and rifampin resistance are tightly linked, with 78% of the reported rifampin-resistant TB cases in 2019 being resistant to isoniazid as well. Rifampin-resistance is primarily due to resistance-conferring mutations in the rifampin-resistance determining region (RRDR) within the rpoB gene. The most frequently observed mutations of the codons in RRDR are 531, 526 and 516. However, alternative more elusive resistance-conferring mutations have been detected. Isoniazid function occurs through the inhibition of mycolic acid synthesis through the NADH-dependent enoyl-acyl carrier protein (ACP)-reductase. This is encoded by the inhA gene. As a result, isoniazid resistance is primarily due to mutations within inhA and the KatG gene or its promoter region - a catalase peroxidase which is required to activate Isoniazid. As MDR in M. tuberculosis becomes increasingly common, the emergence of pre-extensively drug resistant (pre-XDR) and extensively drug resistant (XDR-) TB threatens to exacerbate the public health crisis. XDR-TB is characterised by resistance to both rifampin and Isoniazid, as well second-line fluoroquinolones and at least one additional front-line drug. Thus, the development of alternative therapeutic measures is of utmost priority.
An intrinsic contributor to the antibiotic resistant nature of M. tuberculosis is its unique cell wall. Saturated with long-chain fatty acids or mycolic acids, the mycobacterial cell presents a robust, relatively insoluble barrier. This has led to its synthesis being the target of many antibiotics - such as Isoniazid. However, resistance has emerged to the majority of them. A novel, promising therapeutic target is mycobacterial membrane protein large 3 (MmpL3). The mycobacterial membrane protein large (MmpL) proteins are transmembrane proteins which play a key role in the synthesis of the cell wall and the transport of the associated lipids. Of these, MmpL3 is essential; knock-out of which has been shown to be bactericidal. Due to its essential nature, MmpL3 inhibitors show promise as alternative therapeutic measures in the age of antibiotic resistance. Inhibition of MmpL3 function showed an inability to transport trehalose monomycolate - an essential cell wall lipid - across the plasma membrane. The recently reported structure of MmpL3 revealed resistance-conferring mutations to associate primarily with the transmembrane domain. Although resistance to pre-clinical MmpL3 inhibitors has been detected, analysis of the widespread mutational landscape revealed a low level of environmental resistance. This suggests that MmpL3 inhibitors currently undergoing clinical trials would face little resistance if made available. Additionally, the ability of many MmpL3 inhibitors to work synergistically with other antitubercular drugs presents a ray of hope in combatting the TB crisis.
Host genetics
The nature of the host-pathogen interaction between humans and M. tuberculosis is considered to have a genetic component. A group of rare disorders called Mendelian susceptibility to mycobacterial diseases was observed in a subset of individuals with a genetic defect that results in increased susceptibility to mycobacterial infection.
Early case and twin studies have indicated that genetic components are important in host susceptibility to M. tuberculosis. Recent genome-wide association studies (GWAS) have identified three genetic risk loci, including at positions 11p13 and 18q11. As is common in GWAS, the variants discovered have moderate effect sizes.
DNA repair
As an intracellular pathogen, M. tuberculosis is exposed to a variety of DNA-damaging assaults, primarily from host-generated antimicrobial toxic radicals. Exposure to reactive oxygen species and/or reactive nitrogen species causes different types of DNA damage including oxidation, depurination, methylation, and deamination that can give rise to single- and double-strand breaks (DSBs).
DnaE2 polymerase is upregulated in M. tuberculosis by several DNA-damaging agents, as well as during infection of mice. Loss of this DNA polymerase reduces the virulence of M. tuberculosis in mice. DnaE2 is an error-prone DNA repair polymerase that appears to contribute to M. tuberculosis survival during infection.
The two major pathways employed in repair of DSBs are homologous recombinational repair (HR) and nonhomologous end joining (NHEJ). Macrophage-internalized M. tuberculosis is able to persist if either of these pathways is defective, but is attenuated when both pathways are defective. This indicates that intracellular exposure of M. tuberculosis to reactive oxygen and/or reactive nitrogen species results in the formation of DSBs that are repaired by HR or NHEJ. However deficiency of DSB repair does not appear to impair M. tuberculosis virulence in animal models.
History
M. tuberculosis, then known as the "tubercle bacillus", was first described on 24 March 1882 by Robert Koch, who subsequently received the Nobel Prize in Physiology or Medicine for this discovery in 1905; the bacterium is also known as "Koch's bacillus".
M. tuberculosis has existed throughout history, but the name has changed frequently over time. In 1720, though, the history of tuberculosis started to take shape into what is known of it today; as the physician Benjamin Marten described in his A Theory of Consumption, tuberculosis may be caused by small living creatures transmitted through the air to other patients.
Vaccine
The BCG vaccine (bacille Calmette-Guerin), which was derived from M. bovis, while effective against childhood and severe forms of tuberculosis, has limited success in preventing the most common form of the disease today, adult pulmonary tuberculosis. Because of this, it is primarily used in high tuberculosis incidence regions, and is not a recommended vaccine in the United States due to the low risk of infection. To receive this vaccine in the United States, an individual is required to go through a consultation process with an expert in M. tuberculosis and is only given to those who meet the specific criteria.
Administration of the BCG vaccine has been shown to induced so called "trained immunity", which refers to the enhanced response of the innate immune system. Unlike adaptive immunity, trained immunity involves long-lasting changes in innate immune cells like monocytes and macrophages, which become more responsive to infections. These changes occur through epigenetic reprogramming, such as histone modifications, leading to increased production of pro-inflammatory cytokines.
Research indicates there may be a correlation between BCG vaccination and better immune response to COVID-19.
The DNA vaccine can be used alone or in combination with BCG. DNA vaccines have enough potential to be used with TB treatment and reduce the treatment time in future.
| Biology and health sciences | Gram-positive bacteria | Plants |
392196 | https://en.wikipedia.org/wiki/2MASS | 2MASS | The Two Micron All-Sky Survey, or 2MASS, was an astronomical survey of the whole sky in infrared light. It took place between 1997 and 2001, in two different locations: at the U.S. Fred Lawrence Whipple Observatory on Mount Hopkins, Arizona, and at the Cerro Tololo Inter-American Observatory in Chile, each using a 1.3-meter telescope for the Northern and Southern Hemisphere, respectively. It was conducted in the short-wavelength infrared at three distinct frequency bands (J, H, and K) near 2 micrometres, from which the photometric survey with its HgCdTe detectors derives its name.
2MASS produced an astronomical catalog with over 300 million observed objects, including minor planets of the Solar System, brown dwarfs, low-mass stars, nebulae, star clusters and galaxies. In addition, 1 million objects were cataloged in the 2MASS Extended Source Catalog (2MASX). The cataloged objects are designated with a "2MASS" and "2MASX"-prefix respectively.
Catalog
The final data release for 2MASS occurred in 2003, and is served by the Infrared Science Archive. The goals of this survey included:
Detection of galaxies in the "Zone of Avoidance", a strip of sky obscured in visible light by our own galaxy, the Milky Way.
Detection of brown dwarfs. 2MASS discovered a total of 173, including 2MASS 0939-2448, 2MASS 0415-0935, 2M1207, and 2MASS J04414489+2301513.
An extensive survey of low mass stars, the most common type of star both in our own galaxy and others.
Cataloging of all detected stars and galaxies.
Infrared measurements from the 2MASS survey have been particularly effective at unveiling previously undiscovered star clusters.
Numerical descriptions of point sources (stars, planets, asteroids) and extended sources (galaxies, nebulae) were cataloged by automated computer programs to an average limiting magnitude of about 14. More than 300 million point sources and 1 million extended sources were cataloged. In November 2003, a team of scientists announced the discovery of the Canis Major Dwarf Galaxy, at that time the closest known satellite galaxy to the Milky Way, based on analysis of 2MASS stellar data.
The resulting data and images from the survey are currently in the public domain, and may be accessed online for free by anyone. There is also a list of 2MASS science publications with links to free pre-publication copies of the papers.
2MASS is sponsored by the University of Massachusetts Amherst, the Infrared Processing and Analysis Center (IPAC, run by Jet Propulsion Laboratory (JPL) and Caltech), NASA, and the National Science Foundation (NSF).
| Physical sciences | Surveys and Catalogs | Astronomy |
392309 | https://en.wikipedia.org/wiki/Founder%20effect | Founder effect | In population genetics, the founder effect is the loss of genetic variation that occurs when a new population is established by a very small number of individuals from a larger population. It was first fully outlined by Ernst Mayr in 1942, using existing theoretical work by those such as Sewall Wright. As a result of the loss of genetic variation, the new population may be distinctively different, both genotypically and phenotypically, from the parent population from which it is derived. In extreme cases, the founder effect is thought to lead to the speciation and subsequent evolution of new species.
In the figure shown, the original population has nearly equal numbers of blue and red individuals. The three smaller founder populations show that one or the other color may predominate (founder effect), due to random sampling of the original population. A population bottleneck may also cause a founder effect, though it is not strictly a new population.
The founder effect occurs when a small group of migrants—not genetically representative of the population from which they came—establish in a new area. In addition to founder effects, the new population is often very small, so it shows increased sensitivity to genetic drift, an increase in inbreeding, and relatively low genetic variation.
Founder mutation
In genetics, a founder mutation is a mutation that appears in the DNA of one or more individuals which are founders of a distinct population. Founder mutations initiate with changes that occur in the DNA and can be passed down to other generations. Any organism—from a simple virus to something complex like a mammal—whose progeny carry its mutation has the potential to express the founder effect, for instance a goat or a human.
Founder mutations originate in long stretches of DNA on a single chromosome; indeed, the original haplotype is the whole chromosome. As the generations progress, the proportion of the haplotype that is common to all carriers of the mutation is shortened (due to genetic recombination). This shortening allows scientists to roughly estimate the age of the mutation.
General
The founder effect is a type of genetic drift, occurring when a small group in a population splinters off from the original population and forms a new one. The new colony may have less genetic variation than the original population, and through the random sampling of alleles during reproduction of subsequent generations, continue rapidly towards fixation. The homozygosity increase can be calculated as , where equals inbreeding coefficient and equals population size. This consequence of inbreeding makes the colony more vulnerable to extinction.
The per generation loss of heterozygosity can be calculated as , where equals heterozygosity. The population of the founders of the colony can also be calculated if the loss of heterozygosity from the bottleneck is known using the same equation.
When a newly formed colony is small, its founders can strongly affect the population's genetic makeup far into the future. In humans, who have a slow reproduction rate, the population will remain small for many generations, effectively amplifying the drift effect generation after generation until the population reaches a certain size. The post-bottleneck population growth rate can be calculated as , where equals the number of generations, is the growth rate, is the population equilibrium size, is the natural logarithm base, and is the constant , where is the original size of the founding colony.
Alleles which were present but relatively rare in the original population can move to one of two extremes. The most common one is that the allele is soon lost altogether, but the other possibility is that the allele survives and within a few generations has become much more dispersed throughout the population. The new colony can experience an increase in the frequency of recessive alleles, as well, and as a result, an increased number who are homozygous for certain recessive traits. The equation to calculate reccessive allele frequencies is based on Hardy-Wienberg assumptions.
The variation in gene frequency between the original population and colony may also trigger the two groups to diverge significantly over the course of many generations. As the variance, or genetic distance, increases, the two separated populations may become distinctively different, both genetically and phenotypically, although not only genetic drift, but also natural selection, gene flow and mutation all contribute to this divergence. This potential for relatively rapid changes in the colony's gene frequency led most scientists to consider the founder effect (and by extension, genetic drift) a significant driving force in the evolution of new species. Sewall Wright was the first to attach this significance to random drift and small, newly isolated populations with his shifting balance theory of speciation. Following behind Wright, Ernst Mayr created many persuasive models to show that the decline in genetic variation and small population size accompanying the founder effect were critically important for new species to develop. However, much less support for this view is shown today, since the hypothesis has been tested repeatedly through experimental research, and the results have been equivocal at best. Speciation by genetic drift is a specific case of peripatric speciation which in itself occurs in rare instances. It takes place when a random change in genetic frequency of population favours the survival of a few organisms of the species with rare genes which cause reproductive mutation. These surviving organisms then breed among themselves over a long period of time to create a whole new species whose reproductive systems or behaviors are no longer compatible with the original population.
Serial founder effect
Serial founder effects have occurred when populations migrate over long distances. Such long-distance migrations typically involve relatively rapid movements followed by periods of settlement. The populations in each migration carry only a subset of the genetic diversity carried from previous migrations. As a result, genetic differentiation tends to increase with geographic distance as described by the "isolation by distance" model. The migration of humans out of Africa is characterized by serial founder effects. Africa has the highest degree of human genetic diversity of any continent, which is consistent with an African origin of modern humans.
In island ecology
Founder populations are essential to the study of island biogeography and island ecology. A natural "blank slate" is not easily found, but a classic series of studies on founder population effects was done following the catastrophic 1883 eruption of Krakatoa, which erased all life on the island. Another continuing study has been following the biocolonization of Surtsey, Iceland, a new volcanic island that erupted offshore between 1963 and 1967. An earlier event, the Toba eruption in Sumatra about 73,000 years ago, covered some parts of India with of ash, and must have coated the Nicobar Islands and Andaman Islands, much nearer in the ash fallout cone, with life-smothering layers, forcing the restart of their biodiversity.
However, not all founder effect studies are initiated after a natural disaster; some scientists study the reinstatement of a species that became locally extinct or hadn't existed there before. A study has been in place since 1958 studying the wolf/moose interaction on Isle Royale in Lake Superior after those animals naturally migrated there, perhaps on winter ice. Hajji and others, and Hundertmark & Van Daele, studied the current population statuses of past founder effects in Corsican red deer and Alaskan elk, respectively. Corsican red deer are still listed as an endangered species, decades after a severe bottleneck. They inhabit the Tyrrhenian islands and surrounding mainlands currently, and before the bottleneck, but Hajji and others wanted to know how the deer originally got to the islands, and from what parent population or species they were derived. Through molecular analysis, they were able to determine a possible lineage, with red deer from the islands of Corsica and Sardinia being the most related to one another. These results are promising, as the island of Corsica was repopulated with red deer from the Sardinian island after the original Corsican red deer population became extinct, and the deer now inhabiting the island of Corsica are diverging from those inhabiting Sardinia.
Kolbe and others set up a pair of genetically sequenced and morphologically examined lizards on seven small islands to watch each new population's growth and adaptation to its new environment. Specifically, they were looking at the effects on limb length and perch width, both widely varying phenotypic ranges in the parent population. Unfortunately, immigration did occur, but the founder effect and adaptive differentiation, which could eventually lead to peripatric speciation, were statistically and biologically significant between the island populations after a few years. The authors also point out that although adaptive differentiation is significant, the differences between island populations best reflect the differences between founders and their genetic diversity that has been passed down through the generations.
Founder effects can affect complex traits, such as song diversity. In the Common Myna (Acridotheres tristis), the percentage of unique songs within a repertoire and within‐song complexity were significantly lower in birds from founder populations.
It was found by Tarr et al. (1998) that the loss of heterozygosity of the Laysan finch (Telespiza cantans) after founding events on small islands in the Pacific Ocean closely matched theoretical calculations upon examination of microsatellite loci.
Among human populations
Genetic studies of founder effect have concentrated on discovering ancestral and novel genetic diseases caused by founder effect and, to a lesser degree, on ancestry-related founder effects on populations, races, and ancient migrations, as well other aspects. The founder population could be the common ancestry of race or ethnicity or the forced localizations caused by artificial countries inside the larger group of ancestry, hence causing an original founder effect. Race and specific founder effect mutation diseases are found in all races or ethnicities, and country-specific mutation diseases are caused by increasing homozygosity (the existence of same gene on both chromosomes pairs, hence a recessive disease may increase in just few generations). The genetic abnormality will increase incrementally with the decrease of number of isolated populations making tribe-specific diseases (such as Ashkenazi Jews, Amish, and Bedouins) and novel genetic defects.
In recessive diseases, founder populations where underlying levels of genome-wide homozygosity are high due to shared common ancestry, but also for consanguineous populations that will have large genome-wide homozygous regions due to inbreeding. Having a catalog of disease-associated variation in these populations enables rapid, early, and accurate diagnoses that may improve patient outcomes due to informed clinical management and early interventions. Enclosed communities such as Amish communities, Ashkenazi communities, and relatively isolated islands allow scientists to better understand and further discover the mutated genes that cause these rare diseases and allow them to also discover protective genes as well.
Due to various migrations throughout human history, founder effects are somewhat common among humans in different times and places. The French Canadians of Quebec are a classical example of founder population. Over 150 years of French colonization, between 1608 and 1760, an estimated 8,500 pioneers married and left at least one descendant on the territory. Following the takeover of the colony by the British crown in 1760, immigration from France effectively stopped, but descendants of French settlers continued to grow in number mainly due to their high fertility rate. Intermarriage occurred mostly with the deported Acadians and migrants coming from the British Isles. Since the 20th century, immigration in Quebec and mixing of French Canadians involve people from all over the world. While the French Canadians of Quebec today may be partly of other ancestries, the genetic contribution of the original French founders is predominant, explaining about 90% of regional gene pools, while Acadian (descended from other French settlers in eastern Canada) admixtures contributing 4% British and 2% Native American and other groups contributing less.
In humans, founder effects can arise from cultural isolation, and inevitably, endogamy. For example, the Amish populations in the United States exhibit founder effects because they have grown from a very few founders, have not recruited newcomers, and tend to marry within the community. Though still rare, phenomena such as polydactyly (extra fingers and toes, a symptom of a condition such as Weyers acrodental dysostosis or Ellis–Van Creveld syndrome) are more common in Amish communities than in the American population at large. Maple syrup urine disease affects about one out of 180,000 infants in the general population. Due in part to the founder effect, however, the disease has a much higher prevalence in children of Amish, Mennonite, and Jewish descent. Similarly, a high frequency of fumarase deficiency exists among the 10,000 members of the Fundamentalist Church of Jesus Christ of Latter Day Saints, a community which practices both endogamy and polygyny, where an estimated 75–80% of the community are blood relatives of just two men—founders John Y. Barlow and Joseph Smith Jessop. In South Asia, castes like the Gujjars, the Baniyas and the Pattapu Kapu have estimated founder effects about 10 times as strong as those of Finns and Ashkenazi Jews.
In Africa, many members of the Vadoma tribe inherit ectrodactyly, giving them the nickname of the "two-toed tribe".
The island of Pingelap also suffered a population bottleneck in 1775 following a typhoon that had reduced the population to only 20 people. As a result, complete achromatopsia has a current rate of occurrence of roughly 10%, with an additional 30% being carriers of this recessive condition.
Around 1814, a small group of British colonists founded a settlement on Tristan da Cunha, a group of small islands in the Atlantic Ocean, midway between Africa and South America. One of the early colonists apparently carried a rare, recessive allele for retinitis pigmentosa, a progressive form of blindness that afflicts homozygous individuals. As late as 1961, the majority of the genes in the gene pool on Tristan were still derived from 15 original ancestors; as a consequence of the inbreeding, of 232 people tested in 1961, four were suffering from retinitis pigmentosa. This represents a prevalence of 1 in 58, compared with a worldwide prevalence of around 1 in 4,000.
The abnormally high rate of twin births in Cândido Godói could be explained by the founder effect.
On 31 August 2023, researchers reported, based on genetic studies, that a human ancestor population bottleneck (from a possible 100,000 to 1000 individuals) occurred "around 930,000 and 813,000 years ago ... lasted for about 117,000 years and brought human ancestors close to extinction."
| Biology and health sciences | Basics_4 | Biology |
392547 | https://en.wikipedia.org/wiki/Fuel%20oil | Fuel oil | Fuel oil is any of various fractions obtained from the distillation of petroleum (crude oil). Such oils include distillates (the lighter fractions) and residues (the heavier fractions). Fuel oils include heavy fuel oil (bunker fuel), marine fuel oil (MFO), furnace oil (FO), gas oil (gasoil), heating oils (such as home heating oil), diesel fuel, and others.
The term fuel oil generally includes any liquid fuel that is burned in a furnace or boiler to generate heat (heating oils), or used in an engine to generate power (as motor fuels). However, it does not usually include other liquid oils, such as those with a flash point of approximately , or oils burned in cotton- or wool-wick burners. In a stricter sense, fuel oil refers only to the heaviest commercial fuels that crude oil can yield, that is, those fuels heavier than gasoline (petrol) and naphtha.
Fuel oil consists of long-chain hydrocarbons, particularly alkanes, cycloalkanes, and aromatics. Small molecules, such as those in propane, naphtha, gasoline, and kerosene, have relatively low boiling points, and are removed at the start of the fractional distillation process. Heavier petroleum-derived oils like diesel fuel and lubricating oil are much less volatile and distill out more slowly.
Uses
Oil has many uses; it heats homes and businesses and fuels trucks, ships, and some cars. A small amount of electricity is produced by diesel, but it is more polluting and more expensive than natural gas. It is often used as a backup fuel for peaking power plants in case the supply of natural gas is interrupted or as the main fuel for small electrical generators. In Europe, the use of diesel is generally restricted to cars (about 40%), SUVs (about 90%), and trucks and buses (over 99%). The market for home heating using fuel oil has decreased due to the widespread penetration of natural gas as well as heat pumps. However, it is very common in some areas, such as the Northeastern United States.
Residual fuel oil is less useful because it is so viscous that it has to be heated with a special heating system before use and it may contain relatively high amounts of pollutants, particularly sulfur, which forms sulfur dioxide upon combustion. However, its undesirable properties make it very cheap. In fact, it is the cheapest liquid fuel available. Since it requires heating before use, residual fuel oil cannot be used in road vehicles, boats or small ships, as the heating equipment takes up valuable space and makes the vehicle heavier. Heating the oil is also a delicate procedure, which is impractical on small, fast moving vehicles. However, power plants and large ships are able to use residual fuel oil.
Use of residual fuel oil was more common in the past. It powered boilers, railroad steam locomotives, and steamships. Locomotives, however, have become powered by diesel or electric power; steamships are not as common as they were previously due to their higher operating costs (most LNG carriers use steam plants, as "boil-off" gas emitted from the cargo can be used as a fuel source); and most boilers now use heating oil or natural gas. Some industrial boilers still use it and so do some old buildings, including in New York City. In 2011 New York City estimated that the 1% of its buildings that burned fuel oils No. 4 and No. 6 were responsible for 86% of the soot pollution generated by all buildings in the city. New York made the phase out of these fuel grades part of its environmental plan, PlaNYC, because of concerns for the health effects caused by fine particulates, and all buildings using fuel oil No. 6 had been converted to less polluting fuel by the end of 2015.
Residual fuel's use in electrical generation has also decreased. In 1973, residual fuel oil produced 16.8% of the electricity in the US. By 1983, it had fallen to 6.2%, and , electricity production from all forms of petroleum, including diesel and residual fuel, is only 3% of total production. The decline is the result of price competition with natural gas and environmental restrictions on emissions. For power plants, the costs of heating the oil, extra pollution control and additional maintenance required after burning it often outweigh the low cost of the fuel. Burning fuel oil, particularly residual fuel oil, produces uniformly higher carbon dioxide emissions than natural gas.
Heavy fuel oils continue to be used in the boiler "lighting up" facility in many coal-fired power plants. This use is approximately analogous to using kindling to start a fire. Without performing this act it is difficult to begin the large-scale combustion process.
The chief drawback to residual fuel oil is its high initial viscosity, particularly in the case of No. 6 oil, which requires a correctly engineered system for storage, pumping, and burning. Though it is still usually lighter than water (with a specific gravity usually ranging from 0.95 to 1.03) it is much heavier and more viscous than No. 2 oil, kerosene, or gasoline. No. 6 oil must, in fact, be stored at around heated to before it can be easily pumped, and in cooler temperatures it can congeal into a tarry semisolid. The flash point of most blends of No. 6 oil is, incidentally, about . Attempting to pump high-viscosity oil at low temperatures was a frequent cause of damage to fuel lines, furnaces, and related equipment which were often designed for lighter fuels.
For comparison, BS 2869 Class G heavy fuel oil behaves in similar fashion, requiring storage at , pumping at around and finalizing for burning at around .
Most of the facilities which historically burned No. 6 or other residual oils were industrial plants and similar facilities constructed in the early or mid 20th century, or which had switched from coal to oil fuel during the same time period. In either case, residual oil was seen as a good prospect because it was cheap and readily available. Most of these facilities have subsequently been closed and demolished, or have replaced their fuel supplies with a simpler one such as gas or No. 2 oil. The high sulfur content of No. 6 oil—up to 3% by weight in some extreme cases—had a corrosive effect on many heating systems (which were usually designed without adequate corrosion protection in mind), shortening their lifespans and increasing the polluting effects. This was particularly the case in furnaces that were regularly shut down and allowed to go cold, because the internal condensation produced sulfuric acid.
Environmental cleanups at such facilities are frequently complicated by the use of asbestos insulation on the fuel feed lines. No. 6 oil is very persistent, and does not degrade rapidly. Its viscosity and stickiness also make remediation of underground contamination very difficult, since these properties reduce the effectiveness of methods such as air stripping.
When released into water, such as a river or ocean, residual oil tends to break up into patches or tarballs – mixtures of oil and particulate matter such as silt and floating organic matter – rather than form a single slick. An average of about 5-10% of the material will evaporate within hours of the release, primarily the lighter hydrocarbon fractions. The remainder will then often sink to the bottom of the water column.
Health impacts
Because of the low quality of bunker fuel, when burnt it is especially harmful to the health of humans, causing serious illnesses and deaths. Prior to the IMO's 2020 sulfur cap, shipping industry air pollution was estimated to cause around 400,000 premature deaths each year, from lung cancer and cardiovascular disease, as well as 14 million childhood asthma cases each year.
Even after the introduction of cleaner fuel rules in 2020, shipping air pollution is still estimated to account for around 250,000 deaths each year, and around 6.4 million childhood asthma cases each year.
The hardest hit countries by air pollution from ships are China, Japan, the UK, Indonesia, and Germany. In 2015, shipping air pollution killed an estimated 20,520 people in China, 4,019 people in Japan, and 3,192 people in the UK.
According to an ICCT study, countries located on major shipping lanes are particularly exposed, and can see shipping account for a high percentage of overall deaths from transport sector air pollution. In Taiwan, shipping accounts for 70% of all transport-attributable air pollution deaths in 2015, followed by Morocco at 51%, Malaysia and Japan both at 41%, Vietnam at 39%, and the UK at 38%.
As well as commercial shipping, cruise ships also emit large amounts of air pollution, damaging people's health. Up to 2019, it was reported that the ships of the single largest cruise company, Carnival Corporation & plc, emitted ten times more sulfur dioxide than all of Europe's cars combined.
General classification
United States
Although the following trends generally hold true, different organizations may have different numerical specifications for the six fuel grades. The boiling point and carbon chain length of the fuel increases with fuel oil number. Viscosity also increases with number, and the heaviest oil must be heated for it to flow. Price usually decreases as the fuel number increases.
Number 1 fuel oil is a volatile distillate oil intended for vaporizing pot-type burners and high-performance/clean diesel engines. It is the kerosene refinery cut that boils off immediately after the heavy naphtha cut used for gasoline. This fuel is commonly known as diesel no. 1, kerosene, and jet fuel. Former names include: coal oil, stove oil, and range oil.
Number 2 fuel oil is a distillate home heating oil. Trucks and some cars use similar diesel no. 2 with a cetane number limit describing the ignition quality of the fuel. Both are typically obtained from the light gas oil cut. The name gasoil refers to the original use of this fraction in the late 19th and early 20th centuries—the gas oil cut was used as an enriching agent for carbureted water gas manufacture.
Number 3 fuel oil was a distillate oil for burners requiring low-viscosity fuel. ASTM merged this grade into the number 2 specification, and the term has been rarely used since the mid-20th century.
Number 4 fuel oil is a commercial heating oil for burner installations not equipped with preheaters. It may be obtained from the heavy gas oil cut. This fuel is sometimes known by the Navy specification of Bunker A.
Number 5 fuel oil is a residual-type industrial heating oil requiring preheating to for proper atomization at the burners. It may be obtained from the heavy gas oil cut, or it may be a blend of residual oil with enough number 2 oil to adjust viscosity until it can be pumped without preheating. This fuel is sometimes known by the Navy specification of Bunker B.
Number 6 fuel oil is a high-viscosity residual oil requiring preheating to . Residual means the material remaining after the more valuable cuts of crude oil have boiled off. The residue may contain various undesirable impurities, including 2% water and 0.5% mineral oil. This fuel may be known as residual fuel oil (RFO), by the Navy specification of Bunker C, or by the Pacific Specification of PS-400.
United Kingdom
The British Standard BS 2869, Fuel Oils for Agricultural, Domestic and Industrial Engines, specifies the following fuel oil classes:
Class C1 and C2 fuels are kerosene-type fuels. C1 is for use in flueless appliances (e.g. lamps). C2 is for vaporizing or atomizing burners in appliances connected to flues.
Class A2 fuel is suitable for mobile, off-road applications that are required to use a sulfur-free fuel. Class D fuel is similar to Class A2 and is suitable for use in stationary applications, such as domestic, commercial, and industrial heating. The BS 2869 standard permits Class A2 and Class D fuel to contain up to 7% (V/V) biodiesel (fatty acid methyl ester, FAME), provided the FAME content meets the requirements of the BS EN 14214 standard.
Classes E to H are residual oils for atomizing burners serving boilers or, with the exception of Class H, certain types of larger combustion engines. Classes F to H invariably require heating prior to use; Class E fuel may require preheating, depending on ambient conditions.
Russia
Mazut is a residual fuel oil often derived from Russian petroleum sources and is either blended with lighter petroleum fractions or burned directly in specialized boilers and furnaces. It is also used as a petrochemical feedstock. In the Russian practice, though, "mazut" is an umbrella term roughly synonymous with the fuel oil in general, that covers most of the types mentioned above, except US grades 1 and 2/3, for which separate terms exist (kerosene and diesel fuel/solar oil respectively — Russian practice doesn't differentiate between diesel fuel and heating oil). This is further separated in two grades, "naval mazut" being analogous to US grades 4 and 5, and "furnace mazut", a heaviest residual fraction of the crude, almost exactly corresponding to US Number 6 fuel oil and further graded by viscosity and sulfur content.
Maritime fuel classification
In the maritime field another type of classification is used for fuel oils:
MGO (Marine gas oil) - Roughly equivalent to no. 2 fuel oil, made from distillate only
MDO (Marine diesel oil) - Roughly equivalent to no. 3 fuel oil, a blend of heavy gasoil that may contain very small amounts of black refinery feed stocks, but has a low viscosity up to 12 cSt so it need not be heated for use in internal combustion engines
IFO (Intermediate fuel oil) - Roughly equivalent no. 4 fuel oil, a blend of gasoil and heavy fuel oil, with less gasoil than marine diesel oil
HFO (Heavy fuel oil) - Pure or nearly pure residual oil, roughly equivalent to no. 5 and no. 6 fuel oil
NSFO (Navy special fuel oil) - Another name for no. 5 HFO
MFO (Marine fuel oil) - Another name for no. 6 HFO
Marine diesel oil contains some heavy fuel oil, unlike regular diesels.
Standards and classification
CCAI and CII are two indexes which describe the ignition quality of residual fuel oil, and CCAI is especially often calculated for marine fuels. Despite this, marine fuels are still quoted on the international bunker markets with their maximum viscosity (which is set by the ISO 8217 standard - see below) due to the fact that marine engines are designed to use different viscosities of fuel. The unit of viscosity used is the centistoke (cSt) and the fuels most frequently quoted are listed below in order of cost, the least expensive first.
IFO 380 - Intermediate fuel oil with a maximum viscosity of 380 centistokes (<3.5% sulfur)
IFO 180 - Intermediate fuel oil with a maximum viscosity of 180 centistokes (<3.5% sulfur)
LS 380 - Low-sulfur (<1.0%) intermediate fuel oil with a maximum viscosity of 380 centistokes
LS 180 - Low-sulfur (<1.0%) intermediate fuel oil with a maximum viscosity of 180 centistokes
MDO - Marine diesel oil
MGO - Marine gasoil
LSMGO - Low-sulfur (<0.1%) Marine Gas Oil - The fuel is to be used in EU Ports and Anchorages. EU Sulfur directive 2005/33/EC
ULSMGO - Ultra-Low-Sulfur Marine Gas Oil - referred to as Ultra-Low-Sulfur Diesel (sulfur 0.0015% max) in the US and Auto Gas Oil (sulfur 0.001% max) in the EU. Maximum sulfur allowable in US territories and territorial waters (inland, marine, and automotive) and in the EU for inland use.
The density is also an important parameter for fuel oils since marine fuels are purified before use to remove water and dirt from the oil. Since the purifiers use centrifugal force, the oil must have a density which is sufficiently different from water. Older purifiers work with a fuel having a maximum of 991 kg/m3; with modern purifiers it is also possible to purify oil with a density of 1010 kg/m3.
The first British standard for fuel oil came in 1982. The latest standard is ISO 8217 issued in 2017. The ISO standard describe four qualities of distillate fuels and 10 qualities of residual fuels. Over the years the standards have become stricter on environmentally important parameters such as sulfur content. The latest standard also banned the adding of used lubricating oil (ULO).
Some parameters of marine fuel oils according to ISO 8217 (3. ed 2005):
Maximum sulfur content in the open ocean is 0.5% since January 2020. Maximum sulfur content in designated areas is 0.1% since 1 January 2015. Before then it was 1.00%.
The content of aluminum and silicon is limited because those metals are dangerous for the engine. Those elements are present because some components of the fuel are manufactured with Fluid Catalytic Cracking process, which makes use of catalyst containing aluminum and silicon.
The flash point of all fuels used in the engine room should be at least 60 °C. (DMX is used for things like emergency generators and not normally used in the engine room. Gaseous fuels such as LPG/LNG have special class rules applied to the fuel systems.)
Bunker fuel
Bunker fuel or bunker crude is technically any type of fuel oil used aboard water vessels. Its name is derived from coal bunkers, where the fuel was originally stored. In 2019, large ships consumed 213 million metric tons of bunker fuel. The Australian Customs and the Australian Tax Office defines a bunker fuel as the fuel that powers the engine of a ship or aircraft. Bunker A is No. 4 fuel oil, bunker B is No. 5, and bunker C is No. 6. Since No. 6 is the most common, "bunker fuel" is often used as a synonym for No. 6. No. 5 fuel oil is also called Navy Special Fuel Oil (NSFO) or just navy special; No. 5 or 6 are also commonly called heavy fuel oil (HFO) or furnace fuel oil (FFO); the high viscosity requires heating, usually by a recirculated low pressure steam system, before the oil can be pumped from a bunker tank. Bunkers are rarely labeled this way in modern maritime practice.
Since the 1980s the International Organization for Standardization (ISO) has been the accepted standard for marine fuels (bunkers). The standard is listed under number 8217, with recent updates in 2010 and 2017. The latest edition of bunker fuel specification is ISO 8217: 2017. The standard divides fuels into residual and distillate fuels. The most common residual fuels in the shipping industry are RMG and RMK. The differences between the two are mainly the density and viscosity, with RMG generally being delivered at 380 centistokes or less, and RMK at 700 centistokes or less. Ships with more advanced engines can process heavier, more viscous, and thus cheaper, fuel. Governing bodies around the world, e.g., California, European Union, have established Emission Control Areas (ECA) that limit the maximum sulfur of fuels burned in their ports to limit pollution, reducing the percentage of sulfur and other particulates from 4.5% m/m to as little as 0.10% as of 2015 inside an ECA. As of 2013 3.5% continued to be permitted outside an ECA, but the International Maritime Organization has planned to lower the sulfur content requirement outside the ECAs to 0.5% m/m by 2020. This is where Marine Distillate Fuels and other alternatives to use of heavy bunker fuel come into play. They have similar properties to diesel #2, which is used as road diesel around the world. The most common grades used in shipping are DMA and DMB. Greenhouse gas emissions resulting from the use of international bunker fuels are currently included in national inventories.
Heavy fuel oil is still the primary fuel for cruise ships, a tourism sector that is associated with a clean and friendly image. In stark contrast, the exhaust gas emissions - due to HFO's high sulfur content - result in an eco balance significantly worse than that for individual mobility.
Bunkering
The term "bunkering" broadly relates to storage of petroleum products in tanks (among other, disparate meanings). The precise meaning can be further specialized depending on context. Perhaps the most common, more specialized usage refers to the practice and business of refueling ships. Bunkering operations are located at seaports, and they include the storage of bunker (ship) fuels and the provision of the fuel to vessels.
Alternatively "bunkering" may apply to the shipboard logistics of loading fuel and distributing it among available bunkers (on-board fuel tanks).
Finally, in the context of the oil industry in Nigeria, bunkering has come to refer to the illegal diversion of crude oil (often subsequently refined in makeshift facilities into lighter transportation fuels) by the unauthorized cutting of holes into transport pipelines, often by very crude and hazardous means and causing spills.
As of 2018, some 300 million metric tons of fuel oil is used for ship bunkering. On January 1, 2020, regulations set by the International Marine Organization (IMO) all marine shipping vessels will require the use of very low sulfur fuel oil (0.5% Sulfur) or to install exhaust gas scrubber systems to remove the excess sulfur dioxide. The emissions from ships have generally been controlled by the following sulfur caps on any fuel oil used on board: 3.50% on and after 1 January 2012 and 0.50% on and after 1 January 2020.
Further removal of sulfur translates to additional energy and capital costs and can impact fuel price and availability. If priced correctly the excess cheap yet dirty fuel would find its way into other markets, including displacing some onshore energy production in nations with low environmental protection .
Transportation
Fuel oil is transported worldwide by fleets of oil tankers making deliveries to suitably sized strategic ports such as Houston, US; Singapore; Fujairah, United Arab Emirates; Balboa, Panama, Cristobal, Panama; Sakha, Egypt; Algeciras, Spain and Rotterdam, Netherlands. Where a convenient seaport does not exist, inland transport may be achieved with the use of barges. Lighter fuel oils can also be transported through pipelines. The major physical supply chains of Europe are along the Rhine River.
Environmental issues
Emissions from bunker fuel burning in ships contribute to climate change and to air pollution levels in many port cities, especially where the emissions from industry and road traffic have been controlled. The switch of auxiliary engines from heavy fuel oil to diesel oil at berth can result in large emission reductions, especially for SO2 and PM. CO2 emissions from bunker fuels sold are not added to national GHG emissions. For small countries with large international ports, there is an important difference between the emissions in territorial waters and the total emissions of the fuel sold. At the 1997 Third Conference of the Parties in Kyoto, Japan, countries agreed to exempt bunker fuels, and multilateral military operations, from national emissions totals after insistence from the U.S. climate change delegation for such exemptions.
| Technology | Fuel | null |
392778 | https://en.wikipedia.org/wiki/Microcar | Microcar | Microcar is a term often used for the smallest size of cars, with three or four wheels and often an engine smaller than . Specific types of microcars include bubble cars, cycle cars, invacar, quadricycles and voiturettes. Microcars are often covered by separate regulations to normal cars, having relaxed requirements for registration and licensing.
Predecessors
Voiturette is a term used by some small cars and tricycles manufactured from 1895 to 1910.
Cyclecars are a type of small, lightweight and inexpensive car manufactured mainly between 1910 and the late 1920s.
Europe 1940–1970: Microcars
The first cars to be described as microcars (earlier equivalents were called voiturettes or cyclecars) were built in the United Kingdom and Germany following World War II, and remained popular until the 1960s. They were originally called minicars, but later became known as microcars.
France also produced large numbers of similar tiny vehicles called voiturettes, but they were rarely sold abroad.
Characteristics
Microcars have three or four wheels, although most were three-wheelers which, in many countries, meant that they qualified for lower taxes and were licensed as motorcycles. Another common characteristic is an engine displacement of less than , although several cars with engines up to have also been classified as microcars. Often, the engine was originally designed for a motorcycle.
History
Microcars originated in the years following World War II, when motorcycles transport was commonly used. To provide better weather protection, three-wheeled microcars began increasing in popularity in the United Kingdom, where they could be driven using only a motorcycle licence. One of the first microcars was the 1949 Bond Minicar.
Microcars also became popular in Europe. A demand for cheap personal motorised transport emerged, and their greater fuel efficiency meant that microcars became even more significant when fuel prices rose, partly due to the 1956 Suez Crisis.
The microcar boom lasted until the late 1950s, when larger cars regained popularity. The 1959 introduction of the Mini, which provided greater size and performance at an affordable price, contributed to the decline in popularity of microcars. Production of microcars had largely ceased by the end of the 1960s, due to competition from the Mini, Citroën 2CV, Fiat 500 and Renault 4.
Bubble cars
Several microcars of the 1950s and 1960s were nicknamed bubble cars. This was due to the aircraft-style bubble canopies of vehicles such as the Messerschmitt KR175, Messerschmitt KR200 and the FMR Tg500. Other microcars, such as the Isetta, also had a bubble-like appearance.
German manufacturers of bubble cars included former military aircraft manufacturers Messerschmitt and Heinkel. BMW manufactured the Italian Iso Rivolta Isetta under licence, using an engine based on one from one of their own motorcycles.
The United Kingdom had licence-built right-hand-drive versions of the Heinkel Kabine and the Isetta. The British version of the Isetta was built with only one rear wheel, instead of the narrow-tracked pair of wheels in the normal Isetta design, in order to take advantage of the three-wheel vehicle laws in the United Kingdom. There were also indigenous British three-wheeled microcars, including the Peel Trident.
Examples include the Citroën Prototype C, FMR Tg500, Fuldamobil, Heinkel Kabine, Isetta, Messerschmitt KR175, Messerschmitt KR200, Peel P50, Peel Trident, SMZ S-1L, Trojan 200, and Kleinschnittger F125.
Worldwide 1990–present
Recent microcars include the 2001 Aixam 5xx series, Renault Twizy, Citroën Ami, and XEV Yoyo.
Electric-powered microcars which have reached production include the 1987 CityEl, the 1990 Automobiles ERAD Spacia, the 1999 Corbin Sparrow, the 2001 REVAi, the 2005 Commuter Cars Tango, the 2009 Tazzari Zero and the resurrected Peel P50 of 2011 (the original model of 1962 - 65 being petrol powered).
The Smart Fortwo is often called a microcar in the United States; although it requires a regular licence to drive.
Quadricycle legislation
The European Union introduced the quadricycle category in 1992. In several European countries since then, microcars are classified by governments separately from normal cars, sometimes using the same regulations as motorcycles or mopeds. Therefore, compared with normal cars, microcars often have relaxed requirements for registration and licensing, and can be subject to lower taxes and insurance costs.
Junior cars
Junior cars are motorized cars for children, typically copies of real designs. Originally powered either by electric engines or small internal combustion engines, electric engines currently dominate. From the 1926 Baby Bugatti until today, junior cars are often as expensive as a real car and are built to a higher standard than a ride-in toy car. As with the Bugatti, these are frequently sold directly by real car manufacturers such as Porsche and Ferrari. In the 1990s Aston Martin built a half-scale junior car version of the then-new Aston Martin Virage Volante, with a handmade aluminium body, leather interior, and 160-cc Honda engine. It cost as much as a brand new Mercedes-Benz 190E.
Manufacturers include Pocket Classics, the Little Car Company, Eshelman, and Hackney.
Microcar trucks
There are also a variety of microcar trucks, usually of the "forward control" or van style to provide more cargo room. These might be used for local deliveries on narrow streets that are unsuited to larger vehicles. The Piaggio Ape is a three-wheeled example. The Honda Acty is a four-wheeled example.
Microcars by country of origin
| Technology | Motorized road transport | null |
393012 | https://en.wikipedia.org/wiki/Interurban | Interurban | The interurban (or radial railway in Canada) is a type of electric railway, with tram-like electric self-propelled railcars which run within and between cities or towns. The term "interurban" is usually used in North America, with other terms used outside it. They were very prevalent in many parts of the world before the Second World War and were used primarily for passenger travel between cities and their surrounding suburban and rural communities. Interurban as a term encompassed the companies, their infrastructure, their cars that ran on the rails, and their service. In the United States, the early 1900s interurban was a valuable economic institution, when most roads between towns, many town streets were unpaved, and transportation and haulage was by horse-drawn carriages and carts.
The interurban provided reliable transportation, particularly in winter weather, between towns and countryside. In 1915, of interurban railways were operating in the United States and, for a few years, interurban railways, including the numerous manufacturers of cars and equipment, were the fifth-largest industry in the country. But due to preference given to automobiles, by 1930, most interurbans in North America had stopped operating. A few survived into the 1950s.
Outside of the US, other countries built large networks of high-speed electric tramways that survive today. Notable systems exist in the Low Countries, Poland and Japan, where populations are densely packed around large conurbations such as the Randstad, Upper Silesia, Greater Tokyo Area and Keihanshin. Switzerland, particularly, has a large network of mountain narrow-gauge interurban lines.
In addition, since the early 21st century many tram-train lines are being built, especially in France and Germany but also elsewhere in the world. These can be regarded as interurbans since they run on the streets, like trams, when in cities, while out of them they either share existing railway lines or use lines that were abandoned by the railway companies.
Definition
The term "interurban" was coined by Charles L. Henry, a state senator in Indiana. The Latin, , means "between cities". The interurban fit on a continuum between urban street railways and full-fledged railroads. George W. Hilton and John F. Due identified four characteristics of an interurban:
Electric power for propulsion.
Passenger service as the primary business.
Equipment heavier and faster than urban streetcars.
Operation on tracks in city streets, and in rural areas on roadside tracks or private rights-of-way.
The definition of "interurban" is necessarily blurry. Some town streetcar lines evolved into interurban systems by extending streetcar track from town into the countryside to link adjacent towns together and sometimes by the acquisition of a nearby interurban system. Following initial construction, there was a large amount of consolidation of lines. Other interurban lines effectively became light rail systems with no street running whatsoever, or they became primarily freight-hauling railroads because of a progressive loss of their initial passenger service over the years.
In 1905, the United States Census Bureau defined an interurban as "a street railway having more than half its trackage outside municipal limits." It drew a distinction between "interurban" and "suburban" railroads. A suburban system was oriented toward a city center in a single urban area and served commuter traffic. A regular railroad moved riders from one city center to another city center and also moved a substantial amount of freight.
The typical interurban similarly served more than one city, but it served a smaller region and made more frequent stops, and it was oriented to passenger rather than freight service.
History
Emergence
USA
The development of interurbans in the late nineteenth century resulted from the convergence of two trends: improvements in electric traction, and an untapped demand for transportation in rural areas, particularly in the Midwestern United States. The 1880s saw the first successful deployments of electric traction in streetcar systems. Most of these built on the pioneering work of Frank J. Sprague, who developed an improved method for mounting an electric traction motor and using a trolley pole for pickup. Sprague's work led to widespread acceptance of electric traction for streetcar operations and end of horse-drawn trams.
The late nineteenth-century United States witnessed a boom in agriculture which lasted through the First World War, but transportation in rural areas was inadequate. Conventional steam railroads made limited stops, mostly in towns. These were supplemented by horse and buggies and steamboats, both of which were slow and the latter of which were restricted to navigable rivers. The increased capacity and profitability of the city street railroads offered the possibility of extending them into the countryside to reach new markets, even linking to other towns. The first interurban to emerge in the United States was the Newark and Granville Street Railway in Ohio, which opened in 1889. It was not a major success, but others followed. The development of the automobile was then in its infancy, and to many investors interurbans appeared to be the future of local transportation.
From 1900 to 1916, large networks of interurban lines were constructed across the United States, particularly in the states of Indiana, Ohio, Pennsylvania, Illinois, Iowa, Utah, and California. In 1900, of interurban track existed, but by 1916, this had increased to , a seven-fold expansion. At one point in time beginning in 1901, it was possible to travel from Elkhart Lake, Wisconsin, to Little Falls, New York, exclusively by interurban. During this expansion, in the regions where they operated, particularly in Ohio and Indiana, "...they almost destroyed the local passenger service of the steam railroad." To show how exceptionally busy the interurbans radiating from Indianapolis were in 1926, the immense Indianapolis Traction Terminal (nine roof covered tracks and loading platforms) scheduled 500 trains in and out daily and moved 7 million passengers that year. At their peak the interurbans were the fifth-largest industry in the United States.
Europe
In Belgium, a sprawling, nation-wide system of narrow-gauge vicinal tramways have been built by the NMVB / SNCV to provide transport to smaller towns across the country; the first section opened in 1885. These lines were either electrically operated or run with diesel tramcars, included numerous street-running sections, and inter-operated with local tram networks in the larger cities. Similar to Belgium, Netherlands constructed a large network of interurbans in the early 1900s called streektramlijnen.
In Silesia, today Poland, an extensive interurban system was constructed, starting in 1894 with a narrow-gauge line connecting Gliwice with Piekary Śląskie through Zabrze, Chebzie, Chorzów and Bytom, another connected Katowice and Siemianowice. After four years, in 1898, Kramer & Co. was chosen to start electrification on Katowice Rynek (Kattowitz, Ring) - Zawodzie line, after which Schikora & Wolff completed electrification of four additional lines. In 1912, the first short line was built in Katowice. In 1913, a separate standard gauge system connecting Bytom with suburbs and villages west of the town was launched. After World War I and the Silesian Uprisings, in 1922 the region (and the tram network) was divided between newly independent Poland and Germany, and international services appeared (the last one ran until 1937). In 1928 further standard gauge systems were established in Sosnowiec, Będzin and Dąbrowa Górnicza (the so-called Dabrowa Coal Basin - a region adjoining the Upper Silesian Coal Basin). Between 1928 and 1936 most of the original narrow gauge network was converted to standard, which allowed a connection with the new system in Sosnowiec. By 1931, 47,5% of the narrow-gauge network was reconstructed, with of new standard-gauge track built.
A large network of interurbans started developing around Milan in the late 1800s; they were originally drawn by horses and later powered as steam trams. These initial interurban lines were gradually upgraded with electric traction in the early 1900s with some assistance from Thomas Edison. By the 1930s a vast network of interurbans, the Società Trazione Elettrica Lombarda, connected Milan with surrounding towns.
In the first half of the 20th century, an extensive tramway network covered Northern England, centered on South Lancashire and West Yorkshire. At that time, it was possible to travel entirely by tram from Liverpool Pier Head to the village of Summit, outside Rochdale, a distance of , and with a short bus journey across the Pennines, to connect to another tram network that linked Huddersfield, Halifax and Leeds.
Asia
The first interurban railway in Japan is the Hanshin Electric Railway, built to compete with mainline steam trains on the Osaka to Kobe corridor and completed in 1905. As laws of that time did not allow parallel railways to be built, the line was legally defined as a tramway and included street running at the two ends, but was based on American interurbans and operated with large tramcars on mostly private right-of-way. In the same year, the Keihin Express Railway, or Keikyu, completed a section of what is today part of the Keikyū Main Line between Shinagawa, Tokyo and Kanagawa, Yokohama. This line competes with mainline Japanese National Railways on this busy corridor. Predecessors of the Meitetsu opened their first interurban lines in 1912, what today form parts of the Meitetsu Inuyama Line and Tsushima Line. In 1913, the first section of what will become the Keiō Line opened connecting Chōfu to just outside Shinjuku with street running on what is today the Kōshū Kaidō or National Route 20. Kyushu Electric Railroad, predecessor to Nishitetsu opened its first interurban line in 1914 serving Kitakyushu and surrounding areas, taking heavy inspiration from Hanshin Electric Railway.
Diverging fortunes
Decline in North America
The fortunes of the industry in the US and Canada declined during World War I, particularly into the early 1920s. In 1919 President Woodrow Wilson created the Federal Electric Railways Commission to investigate the financial problems of the industry. The commission submitted its final report to the President in 1920. The commission's report focused on financial management problems and external economic pressures on the industry, and recommended against introducing public financing for the interurban industry. One of the commission's consultants, however, published an independent report stating that private ownership of electric railways had been a failure, and only public ownership would keep the interurbans in business.
Many interurbans had been hastily constructed without realistic projections of income and expenses. They were initially financed by issuing stock and selling bonds. The sale of these financial instruments was often local with salesmen going door to door aggressively pushing this new and exciting "it can't fail" form of transportation. But many of those interurbans did fail, and often quickly. They had poor cash flow from the outset and struggled to raise essential further capital. Interurbans were very vulnerable to acts of nature damaging track and bridges, particularly in the Midwestern United States where flooding was common. Receivership was a common fate when the interurban company could not pay its payroll and other debts, so state courts took over and allowed continued operation while suspending the company's obligation to pay interest on its bonds. In addition, the interurban honeymoon period with the municipalities of 1895–1910 was over.
The large and heavy interurbans, some weighing as much as 65 tons, caused damage to city streets which led to endless disputes over who should bear the repair costs. The rise of private automobile traffic in the middle 1920s aggravated such trends. As the interurban companies struggled financially, they faced rising competition from cars and trucks on newly paved streets and highways, while municipalities sought to alleviate traffic congestion by removing interurbans from city streets. Some companies exited the passenger business altogether to focus on freight, while others sought to buttress their finances by selling surplus electricity in local communities. Several interurbans that attempted to exit the rail business altogether ran afoul of state commissions which required that trains remain running "for the public good", even at a loss.
Many financially weak interurbans did not survive the prosperous 1920s, and most others went bankrupt during the Great Depression. A few struggling lines tried combining to form much larger systems in an attempt to gain operating efficiency and a broader customer base. This occurred in Ohio in year 1930 with the long Cincinnati & Lake Erie Railroad (C&LE), and in Indiana with the very widespread Indiana Railroad. Both had limited success up to 1937–1938 and primarily earned growing revenues from freight rather than passengers. The long Sacramento Northern Railway stopped carrying passengers in 1940 but continued hauling freight into the 1960s by using heavy electric locomotives.
Oliver Jensen, author of American Heritage History of Railroads in America, commented that "...the automobile doomed the interurban whose private tax paying tracks could never compete with the highways that a generous government provided for the motorist." William D. Middleton, in the opening of his 1961 book The Interurban Era, wrote:
Interurban business increased during World War II due to fuel oil rationing and large wartime employment. When the war ended in 1945, riders went back to their automobiles, and most of these lines were finally abandoned. Several systems struggled into the 1950s, including the Baltimore and Annapolis Railroad (passenger service ended 1950), Lehigh Valley Transit Company (1951), West Penn Railways (1952), and the Illinois Terminal Railroad (1958). The West Penn was the largest interurban to operate in the east at and had provided Pittsburgh-area coal country towns with hourly transportation since 1888.
By the 1960s only five remaining interurban lines served commuters in three major metropolitan areas: the North Shore Line and the South Shore Line in Chicago, the Philadelphia Suburban Transportation Company, the New York, Susquehanna and Western Railway in northern New Jersey, and the Long Beach Line in Long Beach and Los Angeles, California (this was the last remaining part of the Pacific Electric system). The Long Beach Line was cut in 1961, the North Shore Line in 1963; the Philadelphia Suburban's route 103 and the NYS&W in New Jersey both ended passenger service in 1966. Today, only the South Shore Line, Norristown High Speed Line (SEPTA Route 100), and SEPTA Routes 101/102 remain.
Some former interurban lines retained freight service for up to several decades after the discontinuance of passenger service. Most were converted to diesel operation, although the Sacramento Northern Railway retained electric freight until 1965.
Consolidation in Europe
After World War II, many interurbans in other countries were also cut back. In Belgium, as intercity transport shifted to cars and buses; the large sections of the vicinal tramways were gradually shut down by the 1980s. At their peak in 1945, the mileage of vicinal tramways reached and exceeded the length of the national railway network.
Sprawling tram networks in the Netherlands extended to neighbouring cities. The vast majority of these lines were not electrified and operated with steam and sometimes petrol or diesel tramcars. Many did not survive the 1920s and 30s for the same reasons American interurbans went bust, but those that did were put back into service during the war years, or at least the remaining parts not yet demolished. One of the largest systems, nicknamed the Blue Tram, was run by the Noord-Zuid-Hollandsche Stoomtramweg-Maatschappij and survived until 1961. Another, the RTM (Rotterdamse Tramweg Maatschappij), which ran in the river delta south-west of Rotterdam, survived until early January 1966. Its demise sparked the rail-related heritage movement in the Netherlands in earnest with the founding of the Tramweg Stichting (Tramway Foundation). Many systems, such as the Hague tramway and the Rotterdam tramway, included long interurban extensions which were operated with larger, higher-speed cars. In close parallel to North America, many systems were abandoned from the 1950s after tram companies switched to buses.
Instigated by the oil crisis in the 1970s, the remaining interurban tramways have enjoyed somewhat of a renaissance in the form of the Sneltram, a modern light rail system that uses high floor, metro-style vehicles and could interoperate into metro networks. Various other interurbans in Europe were folded into local municipal tramway or light rail systems. Switzerland retained many of its interurban lines which now operate as tramways, local railways, S-Bahn, or tram-trains. Milan's vast interurban network was progressively closed in the 1970s but parts of it were reused as the outer parts of the Milan Metro.
Evolution in Japan
Development of Japanese interurbans strayed from their American counterparts from the 1920s. The second boom of interurbans occurred as late as the 1920s and 1930s in Japan, with predecessors of the extensive Kintetsu Railway, Hankyu, Nankai Electric Railway and Odakyu Electric Railway networks starting life during this period. These interurbans, built with straighter tracks, electrified at 1500V and operated using larger cars, were built to even higher standards than the Japanese National Railways network at the time. The (former JNR) Hanwa Line was a wartime acquisition from Nankai, operating 'Super Express' trains on the line at an average speed of , a national record at the time. The old Sendai station terminus of the Miyagi Electric Railway (the predecessor of the JR Senseki Line) was situated in a short single-track underground tunnel built in 1925; this was the first stretch of underground railway in all of Asia, predating the Tokyo Metro Ginza Line by two years. Meanwhile, existing interurbans like the Hanshin Electric Railway started to rebuild their street-running lines into grade-separated exclusive rights-of-way.
After the war, interurbans and other private railway companies received large investments and were allowed to compete not only with mainline trains but also with each other, in order to rejuvenate the country's railway infrastructure and cater to the post-war baby boom. The companies continued their policies of improvement they had followed before the war; lines were reconstructed to allow higher speeds, mainline-sized trains were adopted, street-running sections were rebuilt to elevated or underground rights-of-way, and link lines to growing metro systems were built to allow for through operations. Many of these private railway companies started to adopt standards for full-blown heavy rail lines similar to the national rail network, and, like JR commuter routes, are operated as 'metro-style' commuter railways with mainline-sized vehicles and metro-like frequencies of very few minutes. In 1957, the Odakyu Electric Railway introduced the Odakyu 3000 series SE, the first in a line of luxurious tourist Limited Express trains named 'Romancecars'. These units set a narrow-gauge speed record of on its runs to the mountain spa resort of Hakone.
Many private lines were nationalised during the Second World War. The handful that remained in the hands of JNR after the end of the war – including the Hanwa Line, Senseki Line and the Iida Line – remain outliers on the national JR network, with short station distances, (in the case of the Iida Line) lower-grade infrastructure, and independent termini (such as Aobadori Station and the upper level of Tennōji Station).
Today, trackage of the major sixteen private railways, in many places originally designed as American-style interurban railways, has been upgraded beyond recognition into high capacity urban heavy railways. Private railway companies that started out as interurbans such as Tokyu, Seibu, Odakyu, Hankyu and Tobu; rail transportation now tends to form only a small part of their extensive business empires, which often include real estate, hotels and resorts, and tourist attractions. For example, the Keikyu network has changed unrecognizably from its early days, operating Limited Express services at up to to compete with JR trains, and inter-operating with subway and Keisei Electric Railway trains on through runs extending up to ; the trains retain a red livery based on the Pacific Electric's 'Red Cars', true to the company's interurban roots. The Keiō Line did not fully remove the street running section on the Kōshū Kaidō outside of Shinjuku Station until the 1960s, replacing it with an underground section.
Similar to passenger railway conditions in early 1900s America, intense competition still exists today between private railways and mainline railways operated by the Japan Railways Group along highly congested corridors is a hallmark of suburban railway operations in Japan. For example, on the Osaka to Kobe corridor, JR West competes intensely with both Hankyu Kobe Line and Hanshin Main Line trains in terms of speed, convenience and comfort.
However, a number of urban lines in Japan did close as late as the 2000s, with networks in Kitakyushu and Gifu being shut down.
Today
Austria
Between Vienna and Baden bei Wien the Badner Bahn, operates a classic interurban passenger service, in addition to some freight services. Some interurban lines survive today a local railways in Upper Austria are such as the Linzer Lokalbahn, Lokalbahn Vöcklamarkt–Attersee and Lokalbahn Lambach–Vorchdorf-Eggenberg. While others operate as extension of as local city tramways such as the Traunseebahn which is now connected to the Gmunden Tramway.
Belgium
Today, two surviving interurban networks descending from the vicinal tramways exist in Belgium. The famous Belgian Coast Tram, built in 1885, traverses the entire Belgian coastline and, at a length of , which is the longest tram line in the world. The Charleroi Metro is a never fully completed pre-metro network upgraded and developed from the dense vicinal tramway network around the city.
Canada
Similar to the United States, in Canada most passenger interurbans were removed by the 1950s. One example of continuous passenger service still exists today, the Toronto Transit Commission 501 Queen streetcar line. The western segment of the 501 Streetcar operates largely on what was the T&YRR Port Credit Radial Line, a radial line that remains intact through Etobicoke and up to the border of the neighbouring City of Mississauga, unlike other Toronto radial lines which were all abandoned outside of the 1960s boundary of the City of Toronto.
Germany
In Germany various networks have continued to operate. Karlsruhe revitalized the interurban concept into the Karlsruhe model by renovating two local railways Alb Valley Railway, which already had interoperability with local tram trackage, and the Hardt Railway. Other examples include:
Interurban tram routes serving Mannheim, Heidelberg, Weinheim (Route 5) and Bad Dürkheim (Route 4) as part of the Mannheim/Ludwigshafen Tram System.
Various interurbans upgraded as part of the Rhine-Ruhr Stadtbahn system such as Düsseldorf to Krefeld (U76) and Duisburg (U79), Bochum/Gelsenkirchen (Route 302), Mülheim to Essen and Essen to Gelsenkirchen (Route 107).
interurban route between Cologne to Bonn (Route 16)
Italy
Milan operates one remaining interurban tramway to Limbiate with another interurban route to Carate Brianza/Giussano suspended since 2011. These two lines were once part of large network of interurbans surrounding Milan that were gradually closed in the 1970s.
Japan
In Japan, the vast majority of the major sixteen private railways have roots as interurban electric railway lines that were inspired by the US. But instead of demolishing their trackage in the 1930s, many Japanese interurbans companies upgraded their networks to heavy rail standards, becoming today's large private railways. To this day, private railway companies in Japan operate as highly influential business empires with diverse business interests, encompassing department stores, property developments and even tourist resorts. Many Japanese private railway companies compete with each other for passengers, operate department stores at their city termini, develop suburban properties adjacent to stations they own, and run special tourist attractions with admission included in package deals with rail tickets; similar to operations of large interurban companies in the US during their heyday.
While most interurbans in Japan have been upgraded beyond recognition to high-capacity urban railways, a handful have remained relatively untouched, with street running and using 'lighter-rail' stock. To this day they retain a distinct character similar to classic American interurbans.
These include:
The Keihan Keishin Line, operating between Kyoto and Otsu, with through services to the Tozai Line on the Kyoto side and street running on the Otsu side. Originally was entirely street running into Kyoto but was partly replaced by the opening of the Tozai Subway Line which trains through operate into.
The Keihan Ishiyama Sakamoto Line a primary north south line operating in Otsu with some street running sections in the city center where it connects with the Keishin Line.
The Eizan Electric Railway originally an interurban that once through operated into the network but was isolated after the closure of the Kyoto City Tram. To this day operates light, one or two car consists.
The Enoshima Electric Railway, which is a line with elevated and on-street trackage, and operated with light, two-car articulated trains.
The Fukui Railway Fukubu Line, which operates a variety of express and local services using light rail cars acquired from Nagoya Railroad's defunct interurban network in Gifu, over a line with an extended on-street section.
The Kumamoto Electric Railway, which operates ex-Tokyo subway stock on a line which includes a short on-street section.
The Chikuhō Electric Railroad Line still operates articulated tramcars on private right-of-way, a holdover from its former inter-operation with the defunct Nishitetsu Kitakyushu street tram network.
Netherlands
The only surviving interurban line is also the oldest regional tramway in the Netherlands a line from The Hague to Delft. Which opened as horse-tramway in 1866. Nowadays the line operates as Line 1 of The Hague Tramway. Line E, run by Randstadrail, was an interurban line connecting Rotterdam to The Hague and in the past also to Scheveningen. It now interoperates with the Rotterdam Metro.
Poland
A large interurban network called the Silesian Interurbans still exists today connecting the urban areas of the Upper Silesia. It is one of the largest interurban networks in Europe.
In Łódź region, an interurban tram system connects Łódź, Pabianice, Zgierz and Konstantynów Łódzki, and formerly also Ozorków, Lutomiersk, Aleksandrów Łódzki, Rzgów and Tuszyn.
United States
Only three continuously operating passenger interurbans in the US remain with most being abandoned by the 1950s.
The South Shore Line is now owned by the state of Indiana and uses mainline-sized electric multiple units. Its last section of street running, in Michigan City, Indiana, was finally closed in 2022 for conversion to a grade-separated double-track line.
SEPTA operates two former Philadelphia Suburban lines: the Norristown High Speed Line (Route 100) as an interurban heavy rail line, and the Media–Sharon Hill Line (Routes 101 and 102) as a light rail line.
In Skokie, Illinois, the Skokie Valley portion of the North Shore Line from Dempster Street to Howard Street was acquired by the Chicago Transit Authority in 1963 and now runs as the Yellow Line. The Yellow Line initially operated with third rail from Howard Street to the Skokie Shops and switched to overhead wire for the remainder of the journey to Dempster Street, until 2004 when the overhead wire was replaced with third rail.
Several former interurban rights of way have been reused for modern light rail lines, including the A and E lines of the Los Angeles Metro Rail system and one section of the Baltimore Light RailLink system. Several museums and heritage railways, including the Western Railway Museum, Seashore Trolley Museum, Fox River Trolley Museum, and the Illinois Railway Museum operate restored equipment on former interurban lines. The Iowa Traction Railway still operates freight service today using interurban equipment and infrastructure. The River Line in New Jersey is also considered an interurban.
Switzerland
Switzerland operated a huge number of interurbans which today many have been upgraded into a number of different modes with a few remaining interurban features left. Several still have interurban characteristics such as unprotected alignments next to the road right of ways and/or street running. Today former interurban lines have been upgraded to operate as:
Extensions of local tram system such as Bern Tramway Line 6 to Worb and Baselland Transport Line 10 formed from combining two narrow-gauge interurbans.
Local private railway networks such as Aargau Verkehr, Aare Seeland mobil, Transports Publics du Chablais, Rhaetian Railway, Chemins de fer du Jura or as individual railways like the Lausanne–Bercher line and Waldenburg railway. These lines commonly retain some street running sections and unprotected roadside alignments.
Tram-trains such as the Forch railway which is part of the Zürich S-Bahn operating on the Zürich Tram System in the city center or the St. Gallen–Trogen and Frauenfeld–Wil railways as part of the St. Gallen S-Bahn.
Infrastructure and design
Right of way
Interurbans typically ran along or on a public right-of-way. In towns, interurbans ran in the street, sharing track with existing street railroads. While street running limited acquisition costs, it also required sharp turns and made interurban operations susceptible to traffic congestion. Unlike conventional railroads, it was rare for an interurban to construct long unencumbered stretches of private right-of-way. The torque characteristics of electric operation allowed interurbans to operate on steeper grades than conventional steam locomotives.
Trackage
Compared to conventional steam railroad trackage, interurban rail was light and ballasted lightly, if at all. Most interurbans in North America were built to standard gauge (), but there were exceptions. In Europe narrow-gauge interurbans were more common. In Japan the national mainlines were built to narrow gauge however due to influence from for US interurban operations the first interurban companies in Japan built trackage to standard gauge. This remains the case today with Keikyu and Hanshin, forerunners of Japan's interurbans, still using standard gauge today. Later companies regauged or outright built lines to narrow gauge for better interoperability and consistency with the Japanese mainline standards. Interurbans often used the tracks of existing street railways through city and town streets, and if these street railways were not built to standard, the interurbans had to use the non-standard gauges as well or face the expense of building their own separate trackage through urban areas. Some municipalities deliberately mandated non-standard gauges to prevent freight operations on public streets. In Pennsylvania, many interurbans were constructed using the wide "Pennsylvania trolley gauge" of . In Los Angeles, the Pacific Electric Railway, using standard gauge track, and the Los Angeles Railway, the city's streetcar system, using narrow gauge, shared dual-gauge track in downtown Los Angeles with one rail common to each.
Electrification
Most interurban railways in North America were constructed using the same low-voltage 500 to 600 V DC trolley power in use by the street railways to which they connected. This enabled interurban cars to use the same overhead trolley power on town street car tracks with no electrical change on the cars to accommodate a different voltage. However, higher voltages became necessary to reduce power loss on long-distance transmission lines and routes, though substations were established to boost voltage. In 1905 Westinghouse introduced a 6600 V 25 Hz alternating current (AC) system which a number of railroads adopted. This required fewer substations than DC, but came with higher maintenance costs. The necessary on-board 6600 AC voltage reduction plus AC to DC rectification on each powered car to run DC traction motors added to greater car construction expense plus the operational dangers that such on-board high voltages created.
More common were high-voltage DC systems – usually 1200 V DC, introduced in 1908 by Indianapolis & Louisville Traction Company for their Dixie Flyer and Hoosier Flyer services. In the streets, where high-speed service was not feasible, the cars ran at half speed at 600 V or got a voltage changeover device. such as on the Sacramento Northern. A 2400 V DC third-rail system was installed on the Michigan United Railways's Western Division between Kalamazoo and Grand Rapids in 1915, but was abandoned because of the electrocution potential safety hazard. Even 5000 V DC was tested.
Most interurban cars and freight locomotives collected current from an overhead trolley wire. The cars contacted this wire through the use of a trolley pole or a pantograph. Other designs collected current from a third rail. Some interurbans used both: in open country, the third rail was used and in town, a trolley pole was raised. An example of this was the Chicago, Aurora, and Elgin Railroad where a trolley pole was used in both Aurora and Elgin, Illinois. Third rail was cheaper to maintain and more conductive, but it was more expensive to construct initially and it did not eliminate the need for AC transformers, AC transmission lines, and AC/DC conversion systems. In addition, third rail posed a serious danger to trespassers and animals and was difficult to keep clear of ice.
Trains and equipment
Rolling stock
From 1890 to 1910, roughly, American interurban cars were made of wood and often were very large, weighing up to and measuring as long as . These featured the classic arch-window look with truss-rods and cow-catchers. Three of the best known early companies were Jewett, Niles, and Kuhlman, all of Ohio. These interurbans required a two men crew, an operator and a conductor. By 1910, most new interurban cars were constructed of steel, weighing up to . As competition increased for passengers and costs needed to be reduced in the 1920s, interurban companies and manufacturers attempted to reduce car weight and wind resistance in order to reduce power consumption. The new designs also required only a one-man crew with the operator collecting tickets and making change. The trucks were improved to provide a better ride, acceleration, and top speed but with reduced power consumption. Into the 1930s, better quality and lighter steel and aluminum use reduced weight, and cars were redesigned to ride lower in order to reduce wind resistance. Car design peaked in the early 1930s with the light weight Cincinnati Car Company-built Red Devil cars of the Cincinnati and Lake Erie Railroad.
In addition to passenger cars, interurban companies acquired freight locomotives and line maintenance equipment. A "box motor" was a powered car exclusive for freight that looked like a passenger interurban without windows and had wide side doors for loading freight. A freight motor was geared for power rather than speed and could pull up to six freight cars depending upon the load and grades. Freight cars for interurbans tended to be smaller than those for steam railroads, and they had to have special extended couplers to prevent car corner contact at the very tight grinding turns at city street corners. Maintenance equipment included "line cars" with roof platforms for the trolley wire repair crew, snow plows and snow sweepers with rotating brushes, a car for weed control and to maintain track and ballast. In order to save money, many companies constructed these in their shops using retired or semi-wrecked passenger cars for the frame and the traction motor mounted trucks.
Passenger trains
Passenger interurban service grew out of horse-drawn rail cars operating on city streets. As these routes electrified and extended outside of towns interurbans began to compete with steam railroads for intercity traffic. Interurbans offered more frequent service than steam railroads, with headways of up to one hour or even half an hour. Interurbans also made more stops, usually apart. As interurban routes tended to be single-track this led to extensive use of passing sidings. Single interurban cars would operate with a motorman and conductor, although in later years one-man operation was common. In open country, the typical interurban proceeded at . In towns with the middle of the street operation, speeds were slow and dictated by local ordinance. The result was that the average speed of a scheduled trip was low, as much as under .
Freight trains
Many interurbans did substantial freight business. In 1926, the Cincinnati, Hamilton and Dayton Railway moved of freight per month. By 1929, this had risen to per month. During the 1920s freight revenue helped offset the loss of passenger business to automobiles. A typical interurban freight train consisted of a powered box motor pulling one to four freight cars. These often operated at night as local ordinances forbade daytime freight operation on city streets. Interurban freight in the Midwest was so extensive that Indianapolis constructed a very large freight handling warehouse which all of Indianapolis' seven interurbans companies used.
In literature
In Raymond Chandler's short story The Man who liked Dogs, the narrator trails a suspect in the Los Angeles area:
Carolina Street was away off at the edge of the little beach city. The end of it ran into a disused inter-urban right of way, beyond which stretched a waste of Japanese truck farms. There were just two house in the last block ... the rails were rusted in a forest of weeds.
Similarly in Mandarin's Jade:
The Hotel Tremaine was far out of Santa Monica, near the junk yards. An inter-urban right of way split the street in half, and just as I got to the block that would have the number I had looked up, a two-car train came racketing by at forty-five miles an hour, making almost as much noise as a transport plane taking off. I speeded up beside it and passed the block.
In E.L. Doctorow's Ragtime, a character rides on interurban systems from New York to Boston.
Preservation
Numerous museums, heritage railways and societies have preserved equipment:
California State Railroad Museum
Connecticut Trolley Museum
East Troy Electric Railroad Museum
Electric City Trolley Museum
Fox River Trolley Museum
Fraser Valley Historical Railway Society
Halton County Radial Railway
Illinois Railway Museum
National Museum of Transportation
Orange Empire Railway Museum
Oregon Electric Railway Museum
Pennsylvania Trolley Museum
Rockhill Trolley Museum
Seashore Trolley Museum
Shelburne Falls Trolley Museum
Shore Line Trolley Museum
Vancouver Downtown Historic Railway
Western Railway Museum
| Technology | Rail and cable transport | null |
393061 | https://en.wikipedia.org/wiki/Speleothem | Speleothem | A speleothem (; ) is a geological formation made by mineral deposits that accumulate over time in natural caves. Speleothems most commonly form in calcareous caves due to carbonate dissolution reactions. They can take a variety of forms, depending on their depositional history and environment. Their chemical composition, gradual growth, and preservation in caves make them useful paleoclimatic proxies.
Chemical and physical characteristics
More than 300 variations of cave mineral deposits have been identified. The vast majority of speleothems are calcareous, composed of calcium carbonate (CaCO3) minerals (calcite or aragonite). Less commonly, speleothems are made of calcium sulfate (gypsum or mirabilite) or opal. Speleothems of pure calcium carbonate or calcium sulfate are translucent and colorless. The presence of iron oxide or copper provides a reddish brown color. The presence of manganese oxide can create darker colors such as black or dark brown. Speleothems can also be brown due to the presence of mud and silt.
Many factors impact the shape and color of speleothems, including the chemical composition of the rock and water, water seepage rate, water flow direction, cave temperature, cave humidity, air currents, aboveground climate, and aboveground plant cover. Weaker flows and short travel distances form narrower stalagmites, while heavier flow and a greater fall distance tend to form broader ones.
Formation processes
Most cave chemistry involves calcium carbonate (CaCO3) containing rocks such as limestone or dolomite, composed of calcite or aragonite minerals. Carbonate minerals are more soluble in the presence of higher carbon dioxide (CO2) and lower temperatures. Calcareous speleothems form via carbonate dissolution reactions whereby rainwater reacts with soil CO2 to create weakly acidic water via the reaction:
H2O + CO2 → H2CO3
As the acidic water travels through the calcium carbonate bedrock from the surface to the cave ceiling, it dissolves the bedrock via the reaction:
CaCO3 + H2CO3 → Ca2+ + 2 HCO3−
When the solution reaches a cave, the lower pCO2 in the cave drives the precipitation of CaCO3 via the reaction:
Ca2+ + 2 HCO3− → CaCO3 + H2O + CO2
Over time, the accumulation of these precipitates form dripstones (stalagmites, stalactites), and flowstones, two of the major types of speleothems.
Climate proxies
Speleothem transects can provide paleoclimate records similar to those from ice cores or tree rings. Slow geometrical growth and incorporation of radioactive elements enables speleothems to be accurately and precisely dated over much of the late Quaternary by radiocarbon dating and uranium-thorium dating, as long as the cave is a closed system and the speleothem has not undergone recrystallization. Oxygen (δ18O) and carbon (δ13C) stable isotopes are used to track variation in rainfall temperature, precipitation, and vegetation changes over the past ~500,000 years. The Mg/Ca proxy has likewise been used as a moisture indicator, although its reliability as a palaeohygrometer can be affected by cave ventilation during dry seasons. Variations in precipitation alter the width of speleothem rings: closed rings indicates little rainfall, wider spacing indicates heavier rainfall, and denser rings indicate higher moisture. Drip rate counting and trace element analysis of the water drops record short-term climate variations, such as El Niño–Southern Oscillation (ENSO) climate events. Exceptionally, climate proxy data from the early Permian period have been retrieved from speleothems dated to 289 million years ago sourced from infilled caves exposed by quarrying at the Richards Spur locality in Oklahoma.
Types and categories
Speleothems take various forms, depending on whether the water drips, seeps, condenses, flows, or ponds. Many speleothems are named for their resemblance to man-made or natural objects. Types of speleothems include:
Dripstone is calcium carbonate in the form of stalactites or stalagmites
Stalactites are pointed pendants hanging from the cave ceiling, from which they grow
Soda straws are very thin but long stalactites with an elongated cylindrical shape rather than the usual more conical shape of stalactites
Helictites are stalactites that have a central canal with twig-like or spiral projections that appear to defy gravity
Include forms known as ribbon helictites, saws, rods, butterflies, hands, curly-fries, and "clumps of worms"
Chandeliers are complex clusters of ceiling decorations
Ribbon stalactites, or simply "ribbons", are shaped accordingly
Stalagmites are the "ground-up" counterparts of stalactites, often blunt mounds
Broomstick stalagmites are very tall and spindly
Totem pole stalagmites are also tall and shaped like their namesakes
Fried egg stalagmites are small, typically wider than they are tall
Stalagnate results when stalactites and stalagmites meet or when stalactites reach the floor of the cave
Flowstone is sheet like and found on cave floors and walls
Draperies or curtains are thin, wavy sheets of calcite hanging downward
Bacon is a drapery with variously colored bands within the sheet
Rimstone dams, or gours, occur at stream ripples and form barriers that may contain water
Stone waterfall formations simulate frozen cascades
Cave crystals
Dogtooth spar are large calcite crystals often found near seasonal pools
Frostwork is needle-like growths of calcite or aragonite
Moonmilk is white and cheese-like
Anthodites are flower-like clusters of aragonite crystals
Cryogenic calcite crystals are loose grains of calcite found on the floors of caves formed by segregation of solutes during freezing of water.
Speleogens (technically distinct from speleothems) are formations within caves that are created by the removal of bedrock, rather than as secondary deposits. These include:
Pillars
Scallops
Boneyard
Boxwork
Others
Cave popcorn, also known as "coralloids" or "cave coral", are small, knobby clusters of calcite
Cave pearls are the result of water dripping from high above, causing small "seed" crystals to turn over so often that they form into near-perfect spheres of calcium carbonate
Snottites are colonies of predominantly sulfur oxidizing bacteria and have the consistency of "snot", or mucus
Calcite rafts are thin accumulations of calcite that appear on the surface of cave pools
Hells Bells, a particular speleothem found in the El Zapote cenote of Yucatan in the form of submerged, bell-like shapes
Lava tubes contain speleothems composed of sulfates, mirabilite or opal. When the lava cools, precipitation occurs.
Calthemites
The usual definition of speleothem excludes secondary mineral deposits derived from concrete, lime, mortar, or other calcareous material (e.g. limestone and dolomite) outside the cave environment or in artificial caves (e.g. mines, tunnels), which can have similar shapes and forms as speleothems. Such secondary deposits in man-made structures are termed calthemites. Calthemites are often associated with concrete degradation, or due to leaching of lime, mortar, or other calcareous material.
Gallery
| Physical sciences | Caves | Earth science |
393255 | https://en.wikipedia.org/wiki/Mourning%20dove | Mourning dove | The mourning dove (Zenaida macroura) is a member of the dove family, Columbidae. The bird is also known as the American mourning dove, the rain dove, the chueybird, colloquially as the turtle dove, and it was once known as the Carolina pigeon and Carolina turtledove. It is one of the most abundant and widespread North American birds and a popular gamebird, with more than 20 million birds (up to 70 million in some years) shot annually in the U.S., both for sport and meat. Its ability to sustain its population under such pressure is due to its prolific breeding; in warm areas, one pair may raise up to six broods of two young each in a single year. The wings make an unusual whistling sound upon take-off and landing, a form of sonation. The bird is a strong flier, capable of speeds up to .
Mourning doves are light gray and brown and generally muted in color. Males and females are similar in appearance. The species is generally monogamous, with two squabs (young) per brood. Both parents incubate and care for the young. Mourning doves eat almost exclusively seeds, but the young are fed crop milk by their parents.
Taxonomy
In 1731, the English naturalist Mark Catesby described and illustrated the passenger pigeon and the mourning dove on successive pages of his The Natural History of Carolina, Florida and the Bahama Islands. For the passenger pigeon he used the common name "Pigeon of passage" and the scientific Latin Palumbus migratorius; for the mourning dove he used "Turtle of Carolina" and Turtur carolinensis. In 1743, the naturalist George Edwards included the mourning dove with the English name "long-tail'd dove" and the Latin name Columba macroura in his A Natural History of Uncommon Birds. Edwards's pictures of the male and female doves were drawn from live birds that had been shipped to England from the West Indies. In 1758, when the Swedish naturalist Carl Linnaeus updated his Systema Naturae for the tenth edition, he conflated the two species. He used the Latin name Columba macroura introduced by Edwards as the binomial name but included a description mainly based on Catesby. He cited Edwards's description of the mourning dove and Catesby's description of the passenger pigeon. Linnaeus updated his Systema Naturae again in 1766 for the twelfth edition. He dropped Columba macroura and instead coined Columba migratoria for the passenger pigeon, Columba cariolensis for the mourning dove and Columba marginata for Edwards's mourning dove.
To resolve the confusion over the binomial names of the two species, Francis Hemming proposed in 1952 that the International Commission on Zoological Nomenclature (ICZN) secure the specific name macroura for the mourning dove and migratorius for the passenger pigeon, since this was the intended use by the authors on whose work Linnaeus had based his description. This was accepted by the ICZN, which used its plenary powers to designate the species for the respective names in 1955.
The mourning dove is now placed in the genus Zenaida, introduced in 1838 by the French naturalist Charles Lucien Bonaparte, commemorating his wife Zénaïde. The specific epithet is from the Ancient Greek makros meaning "long" and -ouros meaning "-tailed".
The mourning dove is closely related to the eared dove (Zenaida auriculata) and the Socorro dove (Zenaida graysoni). Some authorities consider them a superspecies, and the three birds are sometimes classified in the separate genus Zenaidura, but the current classification has them as separate species in the genus Zenaida. In addition, the Socorro dove has at times been considered conspecific with the mourning dove, though several differences in behavior, call, and appearance justify separation as two different species. While the three species do form a subgroup of Zenaida, using a separate genus would interfere with the monophyly of Zenaida by making it paraphyletic.
There are five subspecies:
Zenaida macroura marginella (Woodhouse, 1852) – west Canada and west USA to south central Mexico
Zenaida macroura carolinensis (Linnaeus, 1766) – east Canada and east USA, Bermuda, Bahama Islands
Zenaida macroura macroura (Linnaeus, 1758) – (nominate subspecies) Cuba, Hispaniola (Dominican Republic and Haiti), Puerto Rico, Jamaica
Zenaida macroura clarionensis (Townsend, CH, 1890) – Clarion Island (off west Mexico)
Zenaida macroura turturilla (Wetmore, 1956) – Costa Rica, west Panama
The ranges of most of the subspecies overlap a little, with three in the United States or Canada. The West Indian subspecies is found throughout the Greater Antilles. It has recently invaded the Florida Keys. The eastern subspecies is found mainly in eastern North America, as well as Bermuda and the Bahamas. The western subspecies are found in western North America, including parts of Mexico. The Panamanian subspecies is in Central America. The Clarion Island subspecies is found only on Clarion Island, off Mexico's Pacific coast.
The mourning dove is sometimes called the "American mourning dove" to distinguish it from the distantly related mourning collared dove (Streptopelia decipiens) of Africa. It was also formerly known as the "Carolina turtledove" and the "Carolina pigeon". The "mourning" part of its common name comes from its doleful call.
The mourning dove was thought to be the passenger pigeon's closest living relative on morphological grounds until genetic analysis showed Patagioenas pigeons are more closely related. The mourning dove was even suggested to belong to the same genus, Ectopistes, and was listed by some authors as E. carolinensis. The passenger pigeon (Ectopistes migratorius) was hunted to extinction in the early 1900s.
Description
The mourning dove is a medium-sized, slender dove approximately in length. Mourning doves weigh , usually closer to . The mourning dove has a wingspan of 37–45 cm. The elliptical wings are broad, and the head is rounded. Its tail is long and tapered ("macroura" comes from the Greek words for "large" and "tail"). Mourning doves have perching feet, with three toes forward and one reversed. The legs are short and reddish colored. The beak is short and dark, usually a brown-black hue.
The plumage is generally light gray-brown and lighter and pinkish below. The wings have black spotting, and the outer tail feathers are white, contrasting with the black inners. Below the eye is a distinctive crescent-shaped area of dark feathers. The eyes are dark, with light blue skin surrounding them. The adult male has bright purple-pink patches on the neck sides, with light pink coloring reaching the breast. The crown of the adult male is a distinctly bluish-grey color. Females are similar in appearance but with more brown coloring overall and a little smaller than the male. The iridescent feather patches on the neck above the shoulders are nearly absent but can be quite vivid on males. Juvenile birds have a scaly appearance and are generally darker.
Feather colors are generally believed to be relatively static, changing only by small amounts over periods of months. However, a 2011 study argued that since feathers have neither nerves or blood vessels, color changes must be caused by external stimuli. Researchers analyzed how feathers of iridescent mourning doves responded to stimulus changes of adding and evaporating water. As a result, it was discovered that iridescent feather color changed hue, became more chromatic, and increased overall reflectance by almost 50%. Transmission electron microscopy and thin-film models revealed that color is produced by thin-film interference from a single layer of keratin around the edge of feather barbules, under which lies a layer of air and melanosomes. Once the environmental conditions were changed, the most striking morphological difference was a twisting of colored barbules that exposed more of their surface area for reflection, which explains the observed increase in brightness. Overall, the researchers suggest that some plumage colors may be more changeable than previously thought possible.
All five subspecies of the mourning dove look similar and are not easily distinguishable. The nominate subspecies possesses shorter wings and are darker and more buff-colored than the "average" mourning dove. Z. m. carolinensis has longer wings and toes, a shorter beak, and is darker in color. The western subspecies has longer wings, a longer beak and shorter toes, and is more muted and lighter in color. The Panama mourning dove has shorter wings and legs, a longer beak, and is grayer in color. The Clarion Island subspecies possesses larger feet, a larger beak, and is darker brown in color.
Vocalization
This species' call is a distinctive, plaintive , uttered by males to attract females, and it may be mistaken for the call of an owl at first. (Close up, a grating or throat-rattling sound may be heard preceding the first coo.) Other sounds include a nested call () by paired males to attract their female mates to the nest sites, a greeting call (a soft ) by males upon rejoining their mates, and an alarm call (a short ) by either a male or female when threatened. In flight, the wings make a fluttery whistling sound that is hard to hear. The wing whistle is much louder and more noticeable upon take-off and landing. The mourning dove can also 'clap' its wings together when taking off, in a similar manner to the rock dove.
Distribution and habitat
The mourning dove has a large range of nearly . The species is resident throughout the Greater Antilles, most of Mexico, the Continental United States, southern Canada, and the Atlantic archipelago of Bermuda. Much of the Canadian prairie sees these birds in summer only, and southern Central America sees them in winter only. The species is a vagrant in northern Canada, Alaska, and South America. It has been spotted as an accidental at least seven times in the Western Palearctic with records from the British Isles (5), the Azores (1) and Iceland (1). In 1963, the mourning dove was introduced to Hawaii, and in 1998 there was still a small population in North Kona. The mourning dove also appeared on Socorro Island, off the western coast of Mexico, in 1988, sixteen years after the Socorro dove was extirpated from that island.
The mourning dove occupies a wide variety of open and semi-open habitats, such as urban areas, farms, prairie, grassland, and lightly wooded areas. It avoids swamps and thick forest.
Migration
Most mourning doves migrate along flyways over land. Birds in Canada migrate the farthest, probably wintering in Mexico or further south. Those that spend the summer further south are more sedentary, with much shorter migrations. At the southern part of their range, Mourning Doves are present year-round.
Spring migration north runs from March to May. Fall migration south runs from September to November, with immatures moving first, followed by adult females and then by adult males. Migration is usually during the day, in flocks, and at low altitudes.
Behaviour and ecology
Mourning doves sunbathe or rain bathe by lying on the ground or a flat tree limb, leaning over, stretching one wing, and keeping this posture for up to twenty minutes. These birds can also water bathe in shallow pools or birdbaths. Dustbathing is common as well.
Outside the breeding season, mourning doves roost communally in dense deciduous trees or conifers. During sleep, the head rests between the shoulders, close to the body; it is not tucked under the shoulder feathers as in many other species. During the winter in Canada, roosting flights to the roosts in the evening, and out of the roosts in the morning, are delayed on colder days.
Breeding
Courtship begins with a noisy flight by the male, followed by a graceful, circular glide with outstretched wings and head down. After landing, the male will approach the female with a puffed-out breast, bobbing head, and loud calls. Mated pairs will often preen each other's feathers.
The male then leads the female to potential nest sites, and the female will choose one. The female dove builds the nest. The male will fly about, gather material, and bring it to her. The male will stand on the female's back and give the material to the female, who then builds it into the nest. The nest is constructed of twigs, conifer needles, or grass blades, and is of flimsy construction. Mourning doves will sometimes requisition the unused nests of other mourning doves, other birds, or arboreal mammals such as squirrels.
Most nests are in trees, both deciduous and coniferous. Sometimes, they can be found in shrubs, vines, or on artificial constructs like buildings, or hanging flower pots. When there is no suitable elevated object, mourning doves will nest on the ground.
The clutch size is almost always two eggs. Occasionally, however, a female will lay her eggs in the nest of another pair, leading to three or four eggs in the nest. The eggs are white, , long, wide, at laying (5–6% of female body mass). Both sexes incubate, the male from morning to afternoon, and the female the rest of the day and at night. Mourning doves are devoted parents; nests are very rarely left unattended by the adults.
Incubation takes two weeks. The hatched young, called squabs, are strongly altricial, being helpless at hatching and covered with down. Both parents feed the squabs pigeon's milk (crop milk) for the first 3–4 days of life. Thereafter, the crop milk is gradually augmented by seeds. Fledging takes place in about 11–15 days, before the squabs are fully grown but after they are capable of digesting adult food. They stay nearby to be fed by their father for up to two weeks after fledging.
Mourning doves are prolific breeders. In warmer areas, these birds may raise up to six broods in a season. This fast breeding is essential because mortality is high. Each year, mortality can reach 58% a year for adults and 69% for the young.
The mourning dove is generally monogamous and forms strong pair bonds.
Feeding
Mourning doves eat almost exclusively seeds, which make up more than 99% of their diet. Rarely, they will eat snails or insects. Mourning doves generally eat enough to fill their crops and then fly away to digest while resting. They often swallow grit such as fine gravel or sand to assist with digestion. The species usually forages on the ground, walking but not hopping. At bird feeders, mourning doves are attracted to one of the largest ranges of seed types of any North American bird, with a preference for rapeseed, corn, millet, safflower, and sunflower seeds. Mourning doves do not dig or scratch for seeds, though they will push aside ground litter; instead, they eat what is readily visible. They will sometimes perch on plants and eat from there.
Mourning doves show a preference for the seeds of certain species of plant over others. Foods taken in preference to others include pine nuts, sweetgum seeds, and the seeds of pokeberry, amaranth, canary grass, corn, sesame, and wheat. When their favorite foods are absent, mourning doves will eat the seeds of other plants, including buckwheat, rye, goosegrass and smartweed.
Predators and parasites
The primary predators of this species are diurnal birds of prey, such as falcons and hawks. During nesting, corvids, grackles, housecats, or rat snakes will prey on their eggs. Cowbirds rarely parasitize mourning dove nests. Mourning doves reject slightly under a third of cowbird eggs in such nests, and the mourning dove's vegetarian diet is unsuitable for cowbirds.
Mourning doves can be afflicted with several different diseases and parasites, including tapeworms, nematodes, mites, and lice. The mouth-dwelling parasite Trichomonas gallinae is particularly severe. While a mourning dove will sometimes host it without symptoms, it will often cause yellowish growth in the mouth and esophagus that will eventually starve the host to death. Avian pox is a common, insect-vectored disease.
Conservation status
The number of individual mourning doves was estimated to be approximately 475 million in 1994, More recent reports indicate that there were approximately 346 million doves in the US as of September 2023. The large population and its vast range explain why the mourning dove is considered to be of least concern, meaning that the species is not at immediate risk. As a gamebird, the mourning dove is well-managed, with more than 20 million (and up to 40–70 million) shot by hunters each year. However, reporting cautions that mourning doves are in decline in the western United States, and susceptible everywhere in the country due to lead poisoning as they eat spent shot leftover in hunting fields. In some cases, the fields are specifically planted with a favored seed plant to lure them to those sites.
In culture
A Huron/Wyandot legend tells of a maiden named Ayu'ra (modernly spelled Iohara) who used to care for the bird, who came to love her a great deal. One day, she became sick and died. As her spirit traveled across the land to the entrance to the Underworld, all the doves followed her and tried to gain entrance into the Underworld alongside her. Sky Woman, the deity who guards this door, refused them entry, eventually creating smoke to blind them and take Ayu'ra's spirit away without their knowledge. The smoke stained their feathers gray and they have been in mourning for the maiden's loss ever since. The logic behind the story is a play on words—the sound many Native Americans attributed to the bird was "howe howe," and this is also the sound the Iroquoian peoples used to chant over the dead at funerary events.
The eastern mourning dove (Z. m. carolinensis) is Wisconsin's official symbol of peace. The bird is also Michigan's state bird of peace.
The mourning dove appears as the Carolina turtle-dove on plate 286 of Audubon's Birds of America.
| Biology and health sciences | Columbimorphae | Animals |
393258 | https://en.wikipedia.org/wiki/Dirichlet%20series | Dirichlet series | In mathematics, a Dirichlet series is any series of the form
where s is complex, and is a complex sequence. It is a special case of general Dirichlet series.
Dirichlet series play a variety of important roles in analytic number theory. The most usually seen definition of the Riemann zeta function is a Dirichlet series, as are the Dirichlet L-functions. Specifically, the Riemann zeta function ζ(s) is the Dirichlet series of the constant unit function u(n), namely:
where D(u, s) denotes the Dirichlet series of u(n).
It is conjectured that the Selberg class of series obeys the generalized Riemann hypothesis. The series is named in honor of Peter Gustav Lejeune Dirichlet.
Combinatorial importance
Dirichlet series can be used as generating series for counting weighted sets of objects with respect to a weight which is combined multiplicatively when taking Cartesian products.
Suppose that A is a set with a function w: A → N assigning a weight to each of the elements of A, and suppose additionally that the fibre over any natural number under that weight is a finite set. (We call such an arrangement (A,w) a weighted set.) Suppose additionally that an is the number of elements of A with weight n. Then we define the formal Dirichlet generating series for A with respect to w as follows:
Note that if A and B are disjoint subsets of some weighted set (U, w), then the Dirichlet series for their (disjoint) union is equal to the sum of their Dirichlet series:
Moreover, if (A, u) and (B, v) are two weighted sets, and we define a weight function by
for all a in A and b in B, then we have the following decomposition for the Dirichlet series of the Cartesian product:
This follows ultimately from the simple fact that
Examples
The most famous example of a Dirichlet series is
whose analytic continuation to (apart from a simple pole at ) is the Riemann zeta function.
Provided that is real-valued at all natural numbers , the respective real and imaginary parts of the Dirichlet series have known formulas where we write :
Treating these as formal Dirichlet series for the time being in order to be able to ignore matters of convergence, note that we have:
as each natural number has a unique multiplicative decomposition into powers of primes. It is this bit of combinatorics which inspires the Euler product formula.
Another is:
where is the Möbius function. This and many of the following series may be obtained by applying Möbius inversion and Dirichlet convolution to known series. For example, given a Dirichlet character one has
where is a Dirichlet L-function.
If the arithmetic function has a Dirichlet inverse function , i.e., if there exists an inverse function such that the Dirichlet convolution of f with its inverse yields the multiplicative identity
, then the DGF of the inverse function is given by the reciprocal of F:
Other identities include
where is the totient function,
where Jk is the Jordan function, and
where σa(n) is the divisor function. By specialization to the divisor function d = σ0 we have
The logarithm of the zeta function is given by
where Λ(n) is the von Mangoldt function. Similarly, we have that
The logarithmic derivative of the zeta function is then
These last three are special cases of a more general relationship for derivatives of Dirichlet series, given below.
Given the Liouville function λ(n), one has
Yet another example involves Ramanujan's sum:
Another pair of examples involves the Möbius function and the prime omega function:
We have that the Dirichlet series for the prime zeta function, which is the analog to the Riemann zeta function summed only over indices n which are prime, is given by a sum over the Moebius function and the logarithms of the zeta function:
A large tabular catalog listing of other examples of sums corresponding to known Dirichlet series representations is found here.
Examples of Dirichlet series DGFs corresponding to additive (rather than multiplicative) f are given here for the prime omega functions and , which respectively count the number of distinct prime factors of n (with multiplicity or not). For example, the DGF of the first of these functions is expressed as the product of the Riemann zeta function and the prime zeta function for any complex s with :
If f is a multiplicative function such that its DGF F converges absolutely for all , and if p is any prime number, we have that
where is the Moebius function. Another unique Dirichlet series identity generates the summatory function of some arithmetic f evaluated at GCD inputs given by
We also have a formula between the DGFs of two arithmetic functions f and g related by Moebius inversion. In particular, if , then by Moebius inversion we have that . Hence, if F and G are the two respective DGFs of f and g, then we can relate these two DGFs by the formulas:
There is a known formula for the exponential of a Dirichlet series. If is the DGF of some arithmetic f with , then the DGF G is expressed by the sum
where is the Dirichlet inverse of f and where the arithmetic derivative of f is given by the formula for all natural numbers .
Analytic properties
Given a sequence of complex numbers we try to consider the value of
as a function of the complex variable s. In order for this to make sense, we need to consider the convergence properties of the above infinite series:
If is a bounded sequence of complex numbers, then the corresponding Dirichlet series f converges absolutely on the open half-plane Re(s) > 1. In general, if an = O(nk), the series converges absolutely in the half plane Re(s) > k + 1.
If the set of sums
is bounded for n and k ≥ 0, then the above infinite series converges on the open half-plane of s such that Re(s) > 0.
In both cases f is an analytic function on the corresponding open half plane.
In general is the abscissa of convergence of a Dirichlet series if it converges for and diverges for This is the analogue for Dirichlet series of the radius of convergence for power series. The Dirichlet series case is more complicated, though: absolute convergence and uniform convergence may occur in distinct half-planes.
In many cases, the analytic function associated with a Dirichlet series has an analytic extension to a larger domain.
Abscissa of convergence
Suppose
converges for some
Proposition 1.
Proof. Note that:
and define
where
By summation by parts we have
Proposition 2. Define
Then:
is the abscissa of convergence of the Dirichlet series.
Proof. From the definition
so that
which converges as whenever Hence, for every such that diverges, we have and this finishes the proof.
Proposition 3. If converges then as and where it is meromorphic ( has no poles on ).
Proof. Note that
and we have by summation by parts, for
Now find N such that for n > N,
and hence, for every there is a such that for :
Formal Dirichlet series
A formal Dirichlet series over a ring R is associated to a function a from the positive integers to R
with addition and multiplication defined by
where
is the pointwise sum and
is the Dirichlet convolution of a and b.
The formal Dirichlet series form a ring Ω, indeed an R-algebra, with the zero function as additive zero element and the function δ defined by δ(1) = 1, δ(n) = 0 for n > 1 as multiplicative identity. An element of this ring is invertible if a(1) is invertible in R. If R is commutative, so is Ω; if R is an integral domain, so is Ω. The non-zero multiplicative functions form a subgroup of the group of units of Ω.
The ring of formal Dirichlet series over C is isomorphic to a ring of formal power series in countably many variables.
Derivatives
Given
it is possible to show that
assuming the right hand side converges. For a completely multiplicative function ƒ(n), and assuming the series converges for Re(s) > σ0, then one has that
converges for Re(s) > σ0. Here, Λ(n) is the von Mangoldt function.
Products
Suppose
and
If both F(s) and G(s) are absolutely convergent for s > a and s > b then we have
If a = b and ƒ(n) = g(n) we have
Coefficient inversion (integral formula)
For all positive integers , the function f at x, , can be recovered from the Dirichlet generating function (DGF) F of f (or the Dirichlet series over f) using the following integral formula whenever , the abscissa of absolute convergence of the DGF F
It is also possible to invert the Mellin transform of the summatory function of f that defines the DGF F of f to obtain the coefficients of the Dirichlet series (see section below). In this case, we arrive at a complex contour integral formula related to Perron's theorem. Practically speaking, the rates of convergence of the above formula as a function of T are variable, and if the Dirichlet series F is sensitive to sign changes as a slowly converging series, it may require very large T to approximate the coefficients of F using this formula without taking the formal limit.
Another variant of the previous formula stated in Apostol's book provides an integral formula for an alternate sum in the following form for and any real where we denote :
Integral and series transformations
The inverse Mellin transform of a Dirichlet series, divided by s, is given by Perron's formula.
Additionally, if is the (formal) ordinary generating function of the sequence of , then an integral representation for the Dirichlet series of the generating function sequence, , is given by
Another class of related derivative and series-based generating function transformations on the ordinary generating function of a sequence which effectively produces the left-hand-side expansion in the previous equation are respectively defined in.
Relation to power series
The sequence an generated by a Dirichlet series generating function corresponding to:
where ζ(s) is the Riemann zeta function, has the ordinary generating function:
Relation to the summatory function of an arithmetic function via Mellin transforms
If f is an arithmetic function with corresponding DGF F, and the summatory function of f is defined by
then we can express F by the Mellin transform of the summatory function at . Namely, we have that
For and any natural numbers , we also have the approximation to the DGF F of f given by
| Mathematics | Sequences and series | null |
393311 | https://en.wikipedia.org/wiki/Bar%20chart | Bar chart | A bar chart or bar graph is a chart or graph that presents categorical data with rectangular bars with heights or lengths proportional to the values that they represent. The bars can be plotted vertically or horizontally. A vertical bar chart is sometimes called a column chart and has been identified as the prototype of charts.
A bar graph shows comparisons among discrete categories. One axis of the chart shows the specific categories being compared, and the other axis represents a measured value. Some bar graphs present bars clustered in groups of more than one, showing the values of more than one measured variable.
History
Many sources consider William Playfair (1759-1824) to have invented the bar chart and the Exports and Imports of Scotland to and from different parts for one Year from Christmas 1780 to Christmas 1781 graph from his The Commercial and Political Atlas to be the first bar chart in history. Diagrams of the velocity of a constantly accelerating object against time published in The Latitude of Forms (attributed to Jacobus de Sancto Martino or, perhaps, to Nicole Oresme) about 300 years before can be interpreted as "proto bar charts".
Usage
Bar graphs/charts provide a visual presentation of categorical data. Categorical data is a grouping of data into discrete groups, such as months of the year, age group, shoe sizes, and animals. These categories are usually qualitative. In a column (vertical) bar chart, categories appear along the horizontal axis and the height of the bar corresponds to the value of each category.
Bar charts have a discrete domain of categories, and are usually scaled so that all the data can fit on the chart. When there is no natural ordering of the categories being compared, bars on the chart may be arranged in any order. Bar charts arranged from highest to lowest incidence are called Pareto charts.
Grouped (clustered) and stacked
Bar graphs can also be used for more complex comparisons of data with grouped (or "clustered") bar charts, and stacked bar charts.
In grouped (clustered) bar charts, for each categorical group there are two or more bars color-coded to represent a particular grouping. For example, a business owner with two stores might make a grouped bar chart with different colored bars to represent each store: the horizontal axis would show the months of the year and the vertical axis would show revenue.
Alternatively, Stacked bar charts (also known as Composite bar charts) stack bars on top of each other so that the height of the resulting stack shows the combined result. Unlike a grouped bar chart where each factor is displayed next to another, each with their own bar, the stacked bar chart displays multiple data points stacked in a single row or column. This may, for instance, take the form of uniform height bars charting a time series with internal stacked colours indicating the percentage participation of a sub-type of data. Another example would be a time series displaying total numbers, with internal colors indicating participation in the total by sub-types. Stacked bar charts are not suited to data sets having both positive and negative values.
Grouped bar charts usually present the information in the same order in each grouping. Stacked bar charts present the information in the same sequence on each bar.
Variable-width (variwide)
Variable-width bar charts, sometimes abbreviated variwide (bar) charts, are bar charts having bars with non-uniform widths. Generally:
Bars represent quantities with respective rectangles of areas A that are respective arithmetic products of related pairs of
— vertical-axis quantities (A/X) and
— horizontal-axis quantities (X). sajith 300
Anuradhapura 230
Arithmetically, the area of each bar (rectangle) is determined a product of sides' lengths:
(A/X)*X = Area A for each bar
Roles of the vertical and horizontal axes may be reversed, depending on the desired application.
Examples of variable-width bar charts are shown at Wikimedia Commons.
Advantages
Easy to read and interpret: Bar charts are easy to read and interpret, even for people without a background in statistics or data visualization. The bars make it easy to compare values and see trends, making it a useful tool for communicating information to a wide range of audiences.
Can handle large amounts of data: Bar charts can handle large amounts of data and still provide a clear representation of the information. The bars can be made narrow or wide to fit a large number of categories or data points, and the use of color or patterns can make it easier to distinguish between them.
Customizable: Bar charts can be customized to suit the needs of the user. For example, the color, width, and height of the bars can be adjusted to make the chart more visually appealing, and labels and annotations can be added to provide additional information.
Useful for comparing values: Bar charts are particularly useful for comparing values between categories or data points. They allow for quick identification of differences and similarities, making it easy to draw conclusions and make decisions.
Limitations
Limited use for continuous data: Bar charts are not useful for displaying continuous data, such as temperature or time. For continuous data, a line chart or scatter plot may be more appropriate. Bar charts of continuous data with error bars are sometimes referred to as dynamite plots.
Limited use for small sample sizes: Bar charts may not be useful for displaying small sample sizes, as the bars may not accurately represent the data. In such cases, a histogram or box plot may be more appropriate.
May be misleading: Bar charts can be misleading if the scale is not appropriate or if the data is presented in a way that is designed to mislead the viewer. For example, if the y-axis is truncated, the differences between the bars may appear larger than they actually are.
Limited scope for multivariate data: Bar charts can only display one or two variables at a time, making them less useful for displaying multivariate data. In such cases, a scatter plot or heat map may be more appropriate.
| Mathematics | Statistics | null |
16711283 | https://en.wikipedia.org/wiki/BKS%20theory | BKS theory | In the history of quantum mechanics, the Bohr–Kramers–Slater (BKS) theory was perhaps the final attempt at understanding the interaction of matter and electromagnetic radiation on the basis of the so-called old quantum theory, in which quantum phenomena are treated by imposing quantum restrictions on classically describable behaviour. It was advanced in 1924, and sticks to a classical wave description of the electromagnetic field. It was perhaps more a research program than a full physical theory, the ideas that are developed not being worked out in a quantitative way. The purpose of BKS theory was to disprove Einstein's hypothesis of the light quantum.
One aspect, the idea of modelling atomic behaviour under incident electromagnetic radiation using "virtual oscillators" at the absorption and emission frequencies, rather than the (different) apparent frequencies of the Bohr orbits, significantly led Max Born, Werner Heisenberg and Hendrik Kramers to explore mathematics that strongly inspired the subsequent development of matrix mechanics, the first form of modern quantum mechanics. The provocativeness of the theory also generated great discussion and renewed attention to the difficulties in the foundations of the old quantum theory. However, physically the most provocative element of the theory, that momentum and energy would not necessarily be conserved in each interaction but only overall, statistically, was soon shown to be in conflict with experiment.
Walther Bothe won the Nobel Prize in Physics in 1954 for the Bothe–Geiger coincidence experiment that experimentally disproved BKS theory.
Origins
When Albert Einstein introduced the light quantum (photon) in 1905, there was much resistance from the scientific community. However, when in 1923, the Compton effect showed the results could be explained by assuming the light beam behaves as light-quanta and that energy and momentum are conserved, Niels Bohr was still resistant against quantized light, even repudiating it in his 1922 Nobel Prize lecture. So Bohr found a way of using Einstein's approach without also using the light-quantum hypothesis by reinterpreting the principles of energy and momentum conservation as statistical principles. Thus, it was in 1924 that Bohr, Hendrik Kramers and John C. Slater published a provocative description of the interaction of matter and electromagnetic interaction, historically known as the BKS paper that combined quantum transitions and electromagnetic waves with energy and momentum being conserved only on average.
The initial idea of the BKS theory originated with Slater, who proposed to Bohr and Kramers the following elements of a theory of emission and absorption of radiation by atoms, to be developed during his stay in Copenhagen:
Emission and absorption of electromagnetic radiation by matter is realized in agreement with Einstein's photon concept;
A photon emitted by an atom is guided by a classical electromagnetic field (c.f. Louis de Broglie's ideas published September 1923) consisting of spherical waves, thus enabling an explanation of interference;
Even when there are no transitions there exists a classical field to which all atoms contribute; this field contains all frequencies at which an atom can emit or absorb a photon, the probability of such an emission being determined by the amplitude of the corresponding Fourier component of the field; the probabilistic aspect is provisional, to be eliminated when the dynamics of the inside of atoms are better known;
The classical field is not produced by the actual motions of the electrons but by "motions with the frequencies of possible emission and absorption lines" (to be called 'virtual oscillators', creating a field to be referred to as 'virtual' as well).
This fourth point reverts to Max Planck's original view of his quantum introduction in 1900. Planck also did not believe that light was quantized. He believed that a black body had virtual oscillators and that only during interactions between light and the virtual oscillators of the body was the quantum to be considered. Max Planck said in 1911,
Independently, Franz S. Exner had also suggested the statistical validity of energy conservation in the same spirit as the second law of thermodynamics. Erwin Schrödinger, who did his habilitation under the supervision of Exner, was very supportive of the BKS theory. Schrödinger published a paper to provide his own interpretation of the BKS statistical interpretation.
Development with Bohr and Kramers
Slater's main intention seems to have been to reconcile the two conflicting models of radiation, viz. the wave and particle models. He may have had good hopes that his idea with respect to oscillators vibrating at the differences of the frequencies of electron rotations (rather than at the rotation frequencies themselves) might be attractive to Bohr because it solved a problem of the latter's atomic model, even though the physical meaning of these oscillators was far from clear. Nevertheless, Bohr and Kramers had two objections to Slater's proposal:
The assumption that photons exist. Even though Einstein's photon hypothesis could explain in a simple way the photoelectric effect, as well as conservation of energy in processes of de-excitation of an atom followed by excitation of a neighboring one, Bohr had always been reluctant to accept the reality of photons, his main argument being the problem of reconciling the existence of photons with the phenomenon of interference;
The impossibility to account for conservation of energy in a process of de-excitation of an atom followed by excitation of a neighboring one. This impossibility followed from Slater's probabilistic assumption, which did not imply any correlation between processes going on in different atoms.
As Max Jammer puts it, this refocussed the theory "to harmonize the physical picture of the continuous electromagnetic field with the physical picture, not as Slater had proposed of light quanta, but of the discontinuous quantum transitions in the atom." Bohr and Kramers hoped to be able to evade the photon hypothesis on the basis of ongoing work by Kramers to describe "dispersion" (in present-day terms inelastic scattering) of light by means of a classical theory of interaction of radiation and matter. But abandoning the concept of the photon, they instead chose to squarely accept the possibility of non-conservation of energy, and momentum.
Experimental counter-evidence
In the BKS paper the Compton effect was discussed as an application of the idea of "statistical conservation of energy and momentum" in a continuous process of scattering of radiation by a sample of free electrons, where "each of the electrons contributes through the emission of coherent secondary wavelets". Although Arthur Compton had already given an attractive account of his experiment on the basis of the photon picture (including conservation of energy and momentum in individual scattering processes), is it stated in the BKS paper that "it seems at the present state of science hardly justifiable to reject a formal interpretation as that under consideration [i.e. the weaker assumption of statistical conservation] as inadequate". This statement may have prompted experimental physicists to improve `the present state of science' by testing the hypothesis of `statistical energy and momentum conservation'. In any case, already after one year the BKS theory was disproved by coincidence methods studying correlations between the directions into which the emitted radiation and the recoil electron are emitted in individual scattering processes. Such experiments were carried independently, with the Bothe–Geiger coincidence experiment performed by Walther Bothe and Hans Geiger, as well as the experiment by Compton and Alfred W. Simon. They provided experimental evidence pointing in the direction of energy and momentum conservation in individual scattering processes (at least, it was shown that the BKS theory was not able to explain the experimental results). More accurate experiments, performed much later, have also confirmed these results.
Commenting on the experiments, Max von Laue considered that “physics was saved from being led astray.”
From the very beginning, Wolfgang Pauli was extremely critical of the BKS theory, referring to it as the Copenhagen putsch (). In a letter to Kramers, Pauli said that Bohr would have abandoned the theory even if no experiment was ever carried out, arguing that it is the notion of motion and forces that needs to be modified, not the conservation of energy. Pauli could not help to mock the theory, proposing to the Institute of Physics in Copenhague to “fly its flag at half mast on the anniversary of the publication of the work of Bohr, Kramers and Slater.”
As suggested by a letter to Max Born, for Einstein, the corroboration of energy and momentum conservation was probably even more important than his photon hypothesis:
In light of the experimental results, Bohr informed Charles Galton Darwin that "there is nothing else to do than to give our revolutionary efforts as honourable a funeral as possible".
Bohr's reaction, too, was not primarily related to the photon hypothesis. According to Werner Heisenberg, Bohr remarked:
For Bohr the lesson to be learned from the disproof of the BKS theory was not that photons do exist, but rather that the applicability of classical space-time pictures in understanding phenomena within the quantum domain is limited. This theme would become particularly important a few years later in developing the notion of complementarity. According to Heisenberg, Born's statistical interpretation also had its ultimate roots in the BKS theory. Hence, despite its failure the BKS theory still provided an important contribution to the revolutionary transition from classical mechanics to quantum mechanics.
Schrödinger would not abandon the statistical interpretation and would continue to push this theory until the end of his life.
| Physical sciences | Quantum mechanics | Physics |
12274475 | https://en.wikipedia.org/wiki/Variant%20Creutzfeldt%E2%80%93Jakob%20disease | Variant Creutzfeldt–Jakob disease | Variant Creutzfeldt–Jakob disease (vCJD), formerly known as New variant Creutzfeldt–Jakob disease (nvCJD) and referred to colloquially as "mad cow disease" or "human mad cow disease" to distinguish it from its BSE counterpart, is a fatal type of brain disease within the transmissible spongiform encephalopathy family. Initial symptoms include psychiatric problems, behavioral changes, and painful sensations. In the later stages of the illness, patients may exhibit poor coordination, dementia and involuntary movements. The length of time between exposure and the development of symptoms is unclear, but is believed to be years to decades. Average life expectancy following the onset of symptoms is 13 months.
It is caused by prions, which are misfolded proteins. Spread is believed to be primarily due to eating beef infected with bovine spongiform encephalopathy (BSE). Infection is also believed to require a specific genetic susceptibility. Spread may potentially also occur via blood products or contaminated surgical equipment. Diagnosis is by brain biopsy but can be suspected based on certain other criteria. It is different from typical Creutzfeldt–Jakob disease, though both are due to prions.
Treatment for vCJD involves supportive care. As of 2020, 178 cases of vCJD have been recorded in the United Kingdom, due to a 1990s outbreak, and 50 cases in the rest of the world. The disease has become less common since 2000. The typical age of onset is less than 30 years old. It was first identified in 1996 by the National CJD Surveillance Unit in Edinburgh, Scotland.
Signs and symptoms
Initial symptoms include psychiatric problems, behavioral changes, and painful sensations. In the later stages of the illness, patients may exhibit poor coordination, dementia and involuntary movements. The length of time between exposure and the development of symptoms is unclear, but is believed to be years. Average life expectancy following the onset of symptoms is 13 months.
Cause
Tainted beef
In Britain, the primary cause of vCJD has been eating beef tainted with bovine spongiform encephalopathy. A 2012 study by the Health Protection Agency found that around 1 in 2000 had abnormal prions present in appendix cells.
Jonathan Quick, instructor of medicine at the Department of Global Health and Social Medicine at Harvard Medical School, stated that bovine spongiform encephalopathy (BSE) is the first man-made epidemic, or "Frankenstein" disease, because a human decision to feed meat and bone meal to previously herbivorous cattle (as a source of protein) caused what was previously an animal pathogen to enter into the human food chain, and from there to begin causing humans to contract vCJD.
The risk of contracting vCJD from ingestion of cattle products has led to many countries banning the import of beef from countries where BSE has been known to occur, such as the ban on beef from the United States imposed by Japan, South Korea, Mexico, Canada, and other countries in 2003 immediately following the first reported case of BSE in American cattle. Stringent preventative and surveillance practices implemented since then to prevent the disease from entering the human and cattle food chains have caused some to conclude that such bans are unnecessary.
Blood products
As of 2018, evidence suggests that there may be prions in the blood of individuals with vCJD, but this is not the case in individuals with sporadic CJD.
In 2004, a report showed that vCJD can be transmitted by blood transfusions. The finding alarmed healthcare officials because a large epidemic of the disease could result in the near future. A blood test for vCJD infection is possible but is not yet available for screening blood donations. Significant restrictions exist to protect the blood supply. The UK government banned anyone who had received a blood transfusion since January 1980 from donating blood. Since 1999 there has been a ban in the UK for using UK blood to manufacture fractional products such as albumin.
Whilst these restrictions may go some way to preventing a self-sustaining epidemic of secondary infections, the number of infected blood donations is unknown and could be considerable. In June 2013 the government was warned that deaths, then at 176, could rise five-fold through blood transfusions.
On 28 May 2002, the United States Food and Drug Administration instituted a policy that excludes from blood donation anyone having spent at least six months in certain European countries (or three months in the United Kingdom) from 1980 to 1996. Given the large number of U.S. military personnel and their dependents residing in Europe, it was expected that over 7% of donors would be deferred due to the policy. Later changes to this policy first relaxed the restriction to a cumulative total of five years or more of civilian travel in European countries (six months or more if military) then, in 2022, removed it entirely.
In New Zealand, the New Zealand Blood Service (NZBS) in 2000 introduced measures to preclude permanently donors having resided in the United Kingdom (including the Isle of Man and the Channel Islands) for a total of six months or more between January 1980 and December 1996. The measure resulted in ten percent of New Zealand's active blood donors at the time becoming ineligible to donate blood. In 2003, the NZBS further extended restrictions to permanently preclude donors having received a blood transfusion in the United Kingdom since January 1980, and in April 2006, restrictions were further extended to include the Republic of Ireland and France. The restriction was rescinded in late February 2024.
Similar regulations are in place where anyone having spent more than six months for Germany or one year for France living in the UK between January 1980 and December 1996 is permanently banned from donating blood.
In Canada, individuals were formerly ineligible to donate blood or plasma if they had spent a cumulative total of three months or more in the mainland UK or its Crown Dependencies or six months or more in Saudi Arabia from January 1, 1980, through December 31, 1996. They were also ineligible if they had spent a cumulative total of five years or more in France or the Republic of Ireland from January 1, 1980, through 31 December 2001. These restrictions were removed by December 2023.
In Poland, anyone having spent cumulatively six months or longer between 1 January 1980 and 31 December 1996 in the UK, Ireland, or France is permanently barred from donating.
In France, anyone having lived or stayed in the United Kingdom a total of over one year between 1 January 1980 and 31 December 1996 is permanently barred from donating.
In the Czech Republic, anyone having spent more than six months in the UK or France between the years 1980 and 1996 or received transfusion in the UK after the year 1980 is not allowed to donate blood.
In Finland, anyone having lived or stayed in the mainland United Kingdom or its Crown Dependencies for a total of over six months between 1 January 1980 and 31 December 1996 is permanently barred from donating.
Sperm donation
In the U.S., the FDA has banned import of any donor sperm, motivated by a risk of variant Creutzfeldt–Jakob disease, inhibiting the once popular import of Scandinavian sperm. Despite this, the scientific consensus is that the risk is negligible, as there is no evidence Creutzfeldt–Jakob is sexually transmitted.
Occupational contamination
In France, the last two victims of variant Creutzfeldt–Jakob disease, who died in 2019 and 2021, were research technicians at the National Research Institute for Agriculture, Food and the Environment (INRAE). Emilie Jaumain, who died in 2019, at the age of 33, had been the victim of a work accident in 2010, during which she had pricked herself with a tool contaminated with infected brain. The efficacy of this route of contamination has been unambiguously demonstrated in primates. Pierrette C., who died in 2021, had been victim of the same type of work accident. After her diagnosis, a moratorium was initiated in all French laboratories on research activities on infectious prions. In March 2022, INRAE recognized the occupational cause of these two deaths. This raises serious questions about the safety of personnel in these laboratories. Indeed, inspections have noted serious failures in the protection of agents in the face of this deadly risk, and the long incubation period of this disease leads to fears of new cases in the future, hence great concern.
Other causes
Eating other types of brains such as those from squirrels is not recommended as one person contracted vCJD from eating the brain of a squirrel.
Mechanism
Despite the consumption of contaminated beef in the UK being high, vCJD has infected a small number of people. One explanation for this can be found in the genetics of people with the disease. The human PRNP protein which is subverted in prion disease can occur with either methionine or valine at amino acid 129, without any apparent physiological difference. Of the overall white population, about 40% have two methionine-containing alleles, 10% have two valine-containing alleles, and the other 50% are heterozygous at this position. Only a single person with vCJD tested was found to be heterozygous; most of those affected had two copies of the methionine-containing form. It is not yet known whether those unaffected are actually immune or only have a longer incubation period until symptoms appear. Studies in transgenetic mice indicate that all of these genotypes can be affected.
Diagnosis
Definitive
Examination of brain tissue is required to confirm a diagnosis of variant CJD. The following confirmatory features should be present:
Numerous widespread kuru-type amyloid plaques surrounded by vacuoles in both the cerebellum and cerebrum – florid plaques.
Spongiform change and extensive prion protein deposition shown by immunohistochemistry throughout the cerebellum and cerebrum.
Suspected
Current age or age at death less than 55 years (a brain autopsy is recommended, however, for all physician-diagnosed CJD cases).
Psychiatric symptoms at illness onset and/or persistent painful sensory symptoms (frank pain and/or dysesthesia).
Dementia, and development ≥4 months after illness onset of at least two of the following five neurologic signs: poor coordination, myoclonus, chorea, hyperreflexia, or visual signs. (If persistent painful sensory symptoms exist, ≥4 months' delay in the development of the neurologic signs is not required).
A normal or an abnormal EEG, but not the diagnostic EEG changes often seen in classic CJD.
Duration of illness of over 6 months.
Routine investigations do not suggest an alternative, non-CJD diagnosis.
No history of getting human pituitary growth hormone or a dura mater graft from a cadaver.
No history of CJD in a first degree relative or prion protein gene mutation in the person.
Classification
vCJD is a separate condition from classic Creutzfeldt–Jakob disease (though both are caused by PrP prions). Both classic and variant CJD are subtypes of Creutzfeldt–Jakob disease. There are three main categories of CJD disease: sporadic CJD, hereditary CJD, and acquired CJD, with variant CJD being in the acquired group along with iatrogenic CJD. The classic form includes sporadic and hereditary forms. Sporadic CJD is the most common type.
ICD-10 has no separate code for vCJD and such cases are reported under the Creutzfeldt–Jakob disease code (A81.0).
Epidemiology
The Lancet in 2006 suggested that it may take more than 50 years for vCJD to develop, from their studies of kuru, a similar disease in Papua New Guinea. The reasoning behind the claim is that kuru was possibly transmitted through cannibalism in Papua New Guinea when family members would eat the body of a dead relative as a sign of mourning. In the 1950s, cannibalism was banned in Papua New Guinea.
In the late 20th century, however, kuru reached epidemic proportions in certain Papua New Guinean communities, therefore suggesting that vCJD may also have a similar incubation period of 20 to 50 years. A critique to this theory is that while mortuary cannibalism was banned in Papua New Guinea in the 1950s, that does not necessarily mean that the practice ended. Fifteen years later Jared Diamond was informed by Papuans that the practice continued.
These researchers noticed a genetic variation in some people with kuru that has been known to promote long incubation periods. They have also proposed that individuals having contracted CJD in the early 1990s represent a distinct genetic subpopulation, with unusually short incubation periods for bovine spongiform encephalopathy (BSE). This means that there may be many more people with vCJD with longer incubation periods, which may surface many years later.
Prion protein is detectable in lymphoid and appendix tissue up to two years before the onset of neurological symptoms in vCJD. Large scale studies in the UK have yielded an estimated prevalence of 493 per million, higher than the actual number of reported cases. This finding indicates a large number of asymptomatic cases and the need to monitor.
Society and culture
In 1997, a number of people from Kentucky developed CJD. It was discovered that all had consumed squirrel brains. A coincidental relationship between the disease and this dietary practice may have been involved.
In 2008, UK scientists expressed concern over the possibility of a second wave of human cases due to the wide exposure and long incubation of some cases of vCJD. In 2015, a man from New York developed vCJD after eating squirrel brains. From November 2017 to April 2018, four suspected cases of the disease arose in Rochester, NY.
United Kingdom
The first human death from vCJD occurred in the United Kingdom; Wiltshire teenager Stephen Churchill died on 23 May 1995, aged 19. Researchers believe one in 2,000 people in the UK is a carrier of the disease, linked to eating contaminated beef. The survey provides the most robust prevalence measure to date—and identifies abnormal prion protein across a wider age group than found previously and in all genotypes, indicating "infection" may be relatively common. This new study examined over 32,000 anonymous appendix samples. Of these, 16 samples were positive for abnormal prion protein, indicating an overall prevalence of 493 per million population, or one in 2,000 people are likely to be carriers. No difference was seen in different birth cohorts (1941–1960 and 1961–1985), in both sexes, and there was no apparent difference in abnormal prion prevalence in three broad geographical areas. Genetic testing of the 16 positive samples revealed a higher proportion of valine homozygous (VV) genotype on the codon 129 of the gene encoding the prion protein (PRNP) compared with the general UK population. This also differs from the 176 people with vCJD, all of whom to date have been methionine homozygous (MM) genotype. The concern is that individuals with this VV genotype may be susceptible to developing the condition over longer incubation periods.
Human BSE Foundation
In 2000 a voluntary support group was formed by families of people who had died from vCJD. The goal was to support other families going through a similar experience. This support was provided through a National Helpline, a Carer's Guide, a website and a network of family befriending. The support groups had an internet presence at the turn of the 21st century. The driving force behind the foundation was Lester Firkins, whose young son had died from the disease.
In October 2000 the report of the government inquiry into BSE chaired by Lord Phillips was published. The BSE report criticised former Conservative Party Agriculture Ministers John Gummer, John MacGregor and Douglas Hogg. The report concluded that the escalation of BSE into a crisis was the result of intensive farming, particularly with cows being fed with cow and sheep remains. Furthermore, the report was critical of the way the crisis had been handled. There was a reluctance to consider the possibility that BSE could cross the species barrier. The government assured the public that British beef was safe to eat, with agriculture minister John Gummer famously feeding his daughter a burger. The British government were reactive more than proactive in response; the worldwide ban on all British beef exports in March 1996 was a serious economic blow.
The foundation had been calling for compensation to include a care package to help relatives look after those with vCJD. There have been widespread complaints of inadequate health and social services support. Following the Phillips Report in October 2001, the government announced a compensation scheme for British people affected with vCJD. The multi-million-pound financial package was overseen by the vCJD Trust.
A memorial plaque for those who have died due to vCJD was installed in central London in approximately 2000. It is located on the boundary wall of St Thomas' Hospital in Lambeth facing the Riverside Walk of Albert Embankment.
| Biology and health sciences | Prion diseases | Health |
12276398 | https://en.wikipedia.org/wiki/Human%20embryonic%20development | Human embryonic development | Human embryonic development or human embryogenesis is the development and formation of the human embryo. It is characterised by the processes of cell division and cellular differentiation of the embryo that occurs during the early stages of development. In biological terms, the development of the human body entails growth from a one-celled zygote to an adult human being. Fertilization occurs when the sperm cell successfully enters and fuses with an egg cell (ovum). The genetic material of the sperm and egg then combine to form the single cell zygote and the germinal stage of development commences. Human embryonic development covers the first eight weeks of development, which have 23 stages, called Carnegie stages. At the beginning of the ninth week, the embryo is termed a fetus (spelled "foetus" in British English). In comparison to the embryo, the fetus has more recognizable external features and a more complete set of developing organs.
Human embryology is the study of this development during the first eight weeks after fertilization. The normal period of gestation (pregnancy) is about nine months or 36 weeks.
The germinal stage refers to the time from fertilization through the development of the early embryo until implantation is completed in the uterus. The germinal stage takes around 10 days. During this stage, the zygote divides in a process called cleavage. A blastocyst is then formed and implants in the uterus. Embryogenesis continues with the next stage of gastrulation, when the three germ layers of the embryo form in a process called histogenesis, and the processes of neurulation and organogenesis follow.
The entire process of embryogenesis involves coordinated spatial and temporal changes in gene expression, cell growth, and cellular differentiation. A nearly identical process occurs in other species, especially among chordates.
Germinal stage
Fertilization
Fertilization takes place when the spermatozoon has successfully entered the ovum and the two sets of genetic material carried by the gametes fuse together, resulting in the zygote (a single diploid cell). This usually takes place in the ampulla of one of the fallopian tubes. The zygote contains the combined genetic material carried by both the male and female gametes which consists of the 23 chromosomes from the nucleus of the ovum and the 23 chromosomes from the nucleus of the sperm. The 46 chromosomes undergo changes prior to the mitotic division which leads to the formation of the embryo having two cells.
Successful fertilization is enabled by three processes, which also act as controls to ensure species-specificity. The first is that of chemotaxis which directs the movement of the sperm towards the ovum. Secondly, an adhesive compatibility between the sperm and the egg occurs. With the sperm adhered to the ovum, the third process of acrosomal reaction takes place; the front part of the spermatozoan head is capped by an acrosome which contains digestive enzymes to break down the zona pellucida and allow its entry. The entry of the sperm causes calcium to be released which blocks entry to other sperm cells. A parallel reaction takes place in the ovum called the zona reaction. This sees the release of cortical granules that release enzymes which digest sperm receptor proteins, thus preventing polyspermy. The granules also fuse with the plasma membrane and modify the zona pellucida in such a way as to prevent further sperm entry.
Cleavage
The beginning of the cleavage process is marked when the zygote divides through mitosis into two cells. This mitosis continues and the first two cells divide into four cells, then into eight cells and so on. Each division takes from 12 to 24 hours. The zygote is large compared to any other cell and undergoes cleavage without any overall increase in size. This means that with each successive subdivision, the ratio of nuclear to cytoplasmic material increases.
Initially, the dividing cells, called blastomeres ( Greek for sprout), are undifferentiated and aggregated into a sphere enclosed within the zona pellucida of the ovum. When eight blastomeres have formed, they start to compact. They begin to develop gap junctions, enabling them to develop in an integrated way and co-ordinate their response to physiological signals and environmental cues.
When the cells number around sixteen, the solid sphere of cells within the zona pellucida is referred to as a morula.
Blastulation
Cleavage itself is the first stage in blastulation, the process of forming the blastocyst. Cells differentiate into an outer layer of cells called the trophoblast, and an inner cell mass. With further compaction the individual outer blastomeres, the trophoblasts, become indistinguishable. They are still enclosed within the zona pellucida. This compaction serves to make the structure watertight, containing the fluid that the cells will later secrete. The inner mass of cells differentiate to become embryoblasts and polarise at one end. They close together and form gap junctions, which facilitate cellular communication. This polarisation leaves a cavity, the blastocoel, creating a structure that is now termed the blastocyst. (In animals other than mammals, this is called the blastula).
The trophoblasts secrete fluid into the blastocoel. The resulting increase in size of the blastocyst causes it to hatch through the zona pellucida, which then disintegrates. This process is called zona hatching and it takes place on the sixth day of embryo development, immediately before the implantation process. The hatching of the human embryo is supported by proteases secreted by the cells of the blastocyst, which digest proteins of the zona pellucida, giving rise to a hole. Then, due to the rhythmic expansion and contractions of the blastocyst, an increase of the pressure inside the blastocyst itself occurs, the hole expands and finally the blastocyst can emerge from this rigid envelope.
The inner cell mass will give rise to the pre-embryo, the amnion, yolk sac and allantois, while the fetal part of the placenta will form from the outer trophoblast layer. The embryo plus its membranes is called the conceptus, and by this stage the conceptus has reached the uterus. The zona pellucida ultimately disappears completely, and the now exposed cells of the trophoblast allow the blastocyst to attach itself to the endometrium, where it will implant.
The formation of the hypoblast and epiblast, which are the two main layers of the bilaminar germ disc, occurs at the beginning of the second week. Both the embryoblast and the trophoblast will turn into two sub-layers. The inner cells will turn into the hypoblast layer, which will surround the other layer, called the epiblast, and these layers will form the embryonic disc that will develop into the embryo.
The trophoblast will also develop two sub-layers: the cytotrophoblast, which is in front of the syncytiotrophoblast, which in turn lies within the endometrium. Next, another layer called the exocoelomic membrane or Heuser's membrane will appear and surround the cytotrophoblast, as well as the primitive yolk sac. The syncytiotrophoblast will grow and will enter a phase called lacunar stage, in which some vacuoles will appear and be filled by blood in the following days. The development of the yolk sac starts with the hypoblastic flat cells that form the exocoelomic membrane, which will coat the inner part of the cytotrophoblast to form the primitive yolk sac. An erosion of the endothelial lining of the maternal capillaries by the syncytiotrophoblastic cells results in the formation of the maternal sinusoids from where the blood will begin to penetrate and flow into and through the trophoblastic lacunae to give rise to the uteroplacental circulation. Subsequently, new cells derived from yolk sac will be established between trophoblast and exocoelomic membrane and will give rise to extra-embryonic mesoderm, which will form the chorionic cavity.
At the end of the second week of development, some cells of the trophoblast penetrate and form rounded columns into the syncytiotrophoblast. These columns are known as primary villi. At the same time, other migrating cells form into the exocoelomic cavity a new cavity named the secondary or definitive yolk sac, smaller than the primitive yolk sac.
Implantation
After ovulation, the endometrial lining becomes transformed into a secretory lining in preparation of accepting the embryo. It becomes thickened, with its secretory glands becoming elongated, and is increasingly vascular. This lining of the uterine cavity (or womb) is now known as the decidua, and it produces a great number of large decidual cells in its increased interglandular tissue. The blastomeres in the blastocyst are arranged into an outer layer called the trophoblast. The trophoblast then differentiates into an inner layer, the cytotrophoblast, and an outer layer, the syncytiotrophoblast. The cytotrophoblast contains cuboidal epithelial cells and is the source of dividing cells, and the syncytiotrophoblast is a syncytial layer without cell boundaries.
The syncytiotrophoblast implants the blastocyst in the decidual epithelium by projections of chorionic villi, forming the embryonic part of the placenta. The placenta develops once the blastocyst is implanted, connecting the embryo to the uterine wall. The decidua here is termed the decidua basalis; it lies between the blastocyst and the myometrium and forms the maternal part of the placenta. The implantation is assisted by hydrolytic enzymes that erode the epithelium. The syncytiotrophoblast also produces human chorionic gonadotropin, a hormone that stimulates the release of progesterone from the corpus luteum. Progesterone enriches the uterus with a thick lining of blood vessels and capillaries so that it can oxygenate and sustain the developing embryo. The uterus liberates sugar from stored glycogen from its cells to nourish the embryo. The villi begin to branch and contain blood vessels of the embryo. Other villi, called terminal or free villi, exchange nutrients. The embryo is joined to the trophoblastic shell by a narrow connecting stalk that develops into the umbilical cord to attach the placenta to the embryo.
Arteries in the decidua are remodelled to increase the maternal blood flow into the intervillous spaces of the placenta, allowing gas exchange and the transfer of nutrients to the embryo. Waste products from the embryo will diffuse across the placenta.
As the syncytiotrophoblast starts to penetrate the uterine wall, the inner cell mass (embryoblast) also develops. The inner cell mass is the source of embryonic stem cells, which are pluripotent and can develop into any one of the three germ layer cells, and which have the potency to give rise to all the tissues and organs.
Embryonic disc
The embryoblast forms an embryonic disc of two layers, the upper layer is called the epiblast and the lower layer, the hypoblast. The disc is stretched between what will become the amniotic cavity and the yolk sac. The epiblast is adjacent to the trophoblast and made of columnar cells; the hypoblast is closest to the blastocyst cavity and made of cuboidal cells. The epiblast migrates away from the trophoblast downwards, forming the amniotic cavity, the lining of which is formed from amnioblasts developed from the epiblast. The hypoblast is pushed down and forms the yolk sac (exocoelomic cavity) lining. Some hypoblast cells migrate along the inner cytotrophoblast lining of the blastocoel, secreting an extracellular matrix along the way. These hypoblast cells and extracellular matrix are called Heuser's membrane (or the exocoelomic membrane), and they cover the blastocoel to form the yolk sac (or exocoelomic cavity). Cells of the hypoblast migrate along the outer edges of this reticulum and form the extraembryonic mesoderm; this disrupts the extraembryonic reticulum. Soon pockets form in the reticulum, which ultimately coalesce to form the chorionic cavity (extraembryonic coelom).
Gastrulation
The primitive streak, a linear collection of cells formed by the migrating epiblast, appears, and this marks the beginning of gastrulation, which takes place around the seventeenth day (week 3) after fertilization. The process of gastrulation reorganises the two-layer embryo into a three-layer embryo, and also gives the embryo its specific head-to-tail, and front-to-back orientation, by way of the primitive streak which establishes bilateral symmetry. A primitive node (or primitive knot) forms in front of the primitive streak which is the organiser of neurulation. A primitive pit forms as a depression in the centre of the primitive node which connects to the notochord which lies directly underneath. The node has arisen from epiblasts of the amniotic cavity floor, and it is this node that induces the formation of the neural plate which serves as the basis for the nervous system.
The neural plate will form opposite the primitive streak from ectodermal tissue which thickens and flattens into the neural plate. The epiblast in that region moves down into the streak at the location of the primitive pit where the process called ingression, which leads to the formation of the mesoderm takes place. This ingression sees the cells from the epiblast move into the primitive streak in an epithelial-mesenchymal transition; epithelial cells become mesenchymal stem cells, multipotent stromal cells that can differentiate into various cell types. The hypoblast is pushed out of the way and goes on to form the amnion. The epiblast keeps moving and forms a second layer, the mesoderm. The epiblast has now differentiated into the three germ layers of the embryo, so that the bilaminar disc is now a trilaminar disc, the gastrula.
The three germ layers are the ectoderm, mesoderm and endoderm, and are formed as three overlapping flat discs. It is from these three layers that all the structures and organs of the body will be derived through the processes of somitogenesis, histogenesis and organogenesis. The embryonic endoderm is formed by invagination of epiblastic cells that migrate to the hypoblast, while the mesoderm is formed by the cells that develop between the epiblast and endoderm. In general, all germ layers will derive from the epiblast. The upper layer of ectoderm will give rise to the outermost layer of skin, central and peripheral nervous systems, eyes, inner ear, and many connective tissues. The middle layer of mesoderm will give rise to the heart and the beginning of the circulatory system as well as the bones, muscles and kidneys. The inner layer of endoderm will serve as the starting point for the development of the lungs, intestine, thyroid, pancreas and bladder.
Following ingression, a blastopore develops where the cells have ingressed, in one side of the embryo and it deepens to become the archenteron, the first formative stage of the gut. As in all deuterostomes, the blastopore becomes the anus whilst the gut tunnels through the embryo to the other side where the opening becomes the mouth. With a functioning digestive tube, gastrulation is now completed and the next stage of neurulation can begin.
Neurulation
Following gastrulation, the ectoderm gives rise to epithelial and neural tissue, and the gastrula is now referred to as the neurula. The neural plate that has formed as a thickened plate from the ectoderm, continues to broaden and its ends start to fold upwards as neural folds. Neurulation refers to this folding process whereby the neural plate is transformed into the neural tube, and this takes place during the fourth week. They fold, along a shallow neural groove which has formed as a dividing median line in the neural plate. This deepens as the folds continue to gain height, when they will meet and close together at the neural crest. The cells that migrate through the most cranial part of the primitive line form the paraxial mesoderm, which will give rise to the somitomeres that in the process of somitogenesis will differentiate into somites that will form the sclerotomes, the syndetomes, the myotomes and the dermatomes to form cartilage and bone, tendons, dermis (skin), and muscle. The intermediate mesoderm gives rise to the urogenital tract and consists of cells that migrate from the middle region of the primitive line. Other cells migrate through the caudal part of the primitive line and form the lateral mesoderm, and those cells migrating by the most caudal part contribute to the extraembryonic mesoderm.
The embryonic disc begins flat and round, but eventually elongates to have a wider cephalic part and narrow-shaped caudal end. At the beginning, the primitive line extends in cephalic direction and 18 days after fertilization returns caudally until it disappears. In the cephalic portion, the germ layer shows specific differentiation at the beginning of the fourth week, while in the caudal portion it occurs at the end of the fourth week. Cranial and caudal neuropores become progressively smaller until they close completely (by day 26) forming the neural tube.
Development of organs and organ systems
Organogenesis is the development of the organs that begins during the third to eighth week, and continues until birth. Sometimes full development, as in the lungs, continues after birth. Different organs take part in the development of the many organ systems of the body.
Blood
Haematopoietic stem cells that give rise to all the blood cells develop from the mesoderm. The development of blood formation takes place in clusters of blood cells, known as blood islands, in the yolk sac. Blood islands develop outside the embryo, on the umbilical vesicle, allantois, connecting stalk, and chorion, from mesodermal hemangioblasts.
In the centre of a blood island, hemangioblasts form the haematopoietic stem cells that are the precursor to all types of blood cell. In the periphery of a blood island the hemangioblasts differentiate into angioblasts, the precursors to the blood vessels.
Heart and circulatory system
The heart is the first functional organ to develop and starts to beat and pump blood at around 22 days. Cardiac myoblasts and blood islands in the splanchnopleuric mesenchyme on each side of the neural plate give rise to the cardiogenic region.This is a horseshoe-shaped area near to the head of the embryo. By day 19, following cell signalling, two strands begin to form as tubes in this region, as a lumen develops within them. These two endocardial tubes grow and by day 21 have migrated towards each other and fused to form a single primitive heart tube, the tubular heart. This is enabled by the folding of the embryo which pushes the tubes into the thoracic cavity.
Also at the same time that the endocardial tubes are forming, vasculogenesis (the development of the circulatory system) has begun. This starts on day 18 with cells in the splanchnopleuric mesoderm differentiating into angioblasts that develop into flattened endothelial cells. These join to form small vesicles called angiocysts which join up to form long vessels called angioblastic cords. These cords develop into a pervasive network of plexuses in the formation of the vascular network. This network grows by the additional budding and sprouting of new vessels in the process of angiogenesis. Following vasculogenesis and the development of an early vasculature, a stage of vascular remodelling takes place.
The tubular heart quickly forms five distinct regions. From head to tail, these are the infundibulum, bulbus cordis, primitive ventricle, primitive atrium, and the sinus venosus. Initially, all venous blood flows into the sinus venosus, and is propelled from tail to head to the truncus arteriosus. This will divide to form the aorta and pulmonary artery; the bulbus cordis will develop into the right (primitive) ventricle; the primitive ventricle will form the left ventricle; the primitive atrium will become the front parts of the left and right atria and their appendages, and the sinus venosus will develop into the posterior part of the right atrium, the sinoatrial node and the coronary sinus.
Cardiac looping begins to shape the heart as one of the processes of morphogenesis, and this completes by the end of the fourth week. Programmed cell death (apoptosis) at the joining surfaces enables fusion to take place.
In the middle of the fourth week, the sinus venosus receives blood from the three major veins: the vitelline, the umbilical and the common cardinal veins.
During the first two months of development, the interatrial septum begins to form. This septum divides the primitive atrium into a right and a left atrium. Firstly it starts as a crescent-shaped piece of tissue which grows downwards as the septum primum. The crescent shape prevents the complete closure of the atria allowing blood to be shunted from the right to the left atrium through the opening known as the ostium primum. This closes with further development of the system but before it does, a second opening (the ostium secundum) begins to form in the upper atrium enabling the continued shunting of blood.
A second septum (the septum secundum) begins to form to the right of the septum primum. This also leaves a small opening, the foramen ovale which is continuous with the previous opening of the ostium secundum. The septum primum is reduced to a small flap that acts as the valve of the foramen ovale and this remains until its closure at birth. Between the ventricles the septum inferius also forms which develops into the muscular interventricular septum.
Digestive system
The digestive system starts to develop from the third week and by the twelfth week, the organs have correctly positioned themselves.
Respiratory system
The respiratory system develops from the lung bud, which appears in the ventral wall of the foregut about four weeks into development. The lung bud forms the trachea and two lateral growths known as the bronchial buds, which enlarge at the beginning of the fifth week to form the left and right main bronchi. These bronchi in turn form secondary (lobar) bronchi; three on the right and two on the left (reflecting the number of lung lobes). Tertiary bronchi form from secondary bronchi.
While the internal lining of the larynx originates from the lung bud, its cartilages and muscles originate from the fourth and sixth pharyngeal arches.
Urinary system
Kidneys
Three different kidney systems form in the developing embryo: the pronephros, the mesonephros and the metanephros. Only the metanephros develops into the permanent kidney. All three are derived from the intermediate mesoderm.
Pronephros
The pronephros derives from the intermediate mesoderm in the cervical region. It is not functional and degenerates before the end of the fourth week.
Mesonephros
The mesonephros derives from intermediate mesoderm in the upper thoracic to upper lumbar segments. Excretory tubules are formed and enter the mesonephric duct, which ends in the cloaca. The mesonephric duct atrophies in females, but participate in development of the reproductive system in males.
Metanephros
The metanephros appears in the fifth week of development. An outgrowth of the mesonephric duct, the ureteric bud, penetrates metanephric tissue to form the primitive renal pelvis, renal calyces and renal pyramids. The ureter is also formed.
Bladder and urethra
Between the fourth and seventh weeks of development, the urorectal septum divides the cloaca into the urogenital sinus and the anal canal. The upper part of the urogenital sinus forms the bladder, while the lower part forms the urethra.
Reproductive system
Integumentary system
The superficial layer of the skin, the epidermis, is derived from the ectoderm. The deeper layer, the dermis, is derived from mesenchyme.
The formation of the epidermis begins in the second month of development and it acquires its definitive arrangement at the end of the fourth month. The ectoderm divides to form a flat layer of cells on the surface known as the periderm. Further division forms the individual layers of the epidermis.
The mesenchyme that will form the dermis is derived from three sources:
The mesenchyme that forms the dermis in the limbs and body wall derives from the lateral plate mesoderm
The mesenchyme that forms the dermis in the back derives from paraxial mesoderm
The mesenchyme that forms the dermis in the face and neck derives from neural crest cells
Nervous system
Late in the fourth week, the superior part of the neural tube bends ventrally as the cephalic flexure at the level of the future midbrain—the mesencephalon. Above the mesencephalon is the prosencephalon (future forebrain) and beneath it is the rhombencephalon (future hindbrain).
Cranial neural crest cells migrate to the pharyngeal arches as neural stem cells, where they develop in the process of neurogenesis into neurons.
The optical vesicle (which eventually becomes the optic nerve, retina and iris) forms at the basal plate of the prosencephalon. The alar plate of the prosencephalon expands to form the cerebral hemispheres (the telencephalon) whilst its basal plate becomes the diencephalon. Finally, the optic vesicle grows to form an optic outgrowth.
Development of physical features
Face and neck
From the third to the eighth week the face and neck develop.
Ears
The inner ear, middle ear and outer ear have distinct embryological origins.
Inner ear
At about 22 days into development, the ectoderm on each side of the rhombencephalon thickens to form otic placodes. These placodes invaginate to form otic pits, and then otic vesicles. The otic vesicles then form ventral and dorsal components.
The ventral component forms the saccule and the cochlear duct. In the sixth week of development the cochlear duct emerges and penetrates the surrounding mesenchyme, travelling in a spiral shape until it forms 2.5 turns by the end of the eighth week. The saccule is the remaining part of the ventral component. It remains connected to the cochlear duct via the narrow ductus reuniens.
The dorsal component forms the utricle and semicircular canals.
Middle ear
The first pharyngeal pouch lengthens and expands to form the tubotympanic recess. This recess differentiates to form most of the tympanic cavity of the middle ear, and all of the Eustachian or auditory tube. The narrow auditory tube connects the tympanic cavity to the pharynx.
The bones of the middle ear, the ossicles, derive from the cartilages of the pharyngeal arches. The malleus and incus derive from the cartilage of the first pharyngeal arch, whereas the stapes derives from the cartilage of the second pharyngeal arch.
Outer ear
The external auditory meatus develops from the dorsal portion of the first pharyngeal cleft. Six auricular hillocks, which are mesenchymal proliferations at the dorsal aspects of the first and second pharyngeal arches, form the auricle of the ear.
Eyes
The eyes begin to develop from the third week to the tenth week.
Limbs
At the end of the fourth week limb development begins. Limb buds appear on the ventrolateral aspect of the body. They consist of an outer layer of ectoderm and an inner part consisting of mesenchyme which is derived from the parietal layer of lateral plate mesoderm. Ectodermal cells at the distal end of the buds form the apical ectodermal ridge, which creates an area of rapidly proliferating mesenchymal cells known as the progress zone. Cartilage (some of which ultimately becomes bone) and muscle develop from the mesenchyme.
Clinical significance
Toxic exposures in the embryonic period can be the cause of major congenital malformations, since the precursors of the major organ systems are now developing.
Each cell of the preimplantation embryo has the potential to form all of the different cell types in the developing embryo. This cell potency means that some cells can be removed from the preimplantation embryo and the remaining cells will compensate for their absence. This has allowed the development of a technique known as preimplantation genetic diagnosis, whereby a small number of cells from the preimplantation embryo created by IVF, can be removed by biopsy and subjected to genetic diagnosis. This allows embryos that are not affected by defined genetic diseases to be selected and then transferred to the mother's uterus.
Sacrococcygeal teratomas, tumours formed from different types of tissue, that can form, are thought to be related to primitive streak remnants, which ordinarily disappear.
First arch syndromes are congenital disorders of facial deformities, caused by the failure of neural crest cells to migrate to the first pharyngeal arch.
Spina bifida a congenital disorder is the result of the incomplete closure of the neural tube.
Vertically transmitted infections can be passed from the mother to the unborn child at any stage of its development.
Hypoxia a condition of inadequate oxygen supply can be a serious consequence of a preterm or premature birth.
| Biology and health sciences | Human reproduction | Biology |
1686272 | https://en.wikipedia.org/wiki/Chemical%20biology | Chemical biology | Chemical biology is a scientific discipline between the fields of chemistry and biology. The discipline involves the application of chemical techniques, analysis, and often small molecules produced through synthetic chemistry, to the study and manipulation of biological systems. Although often confused with biochemistry, which studies the chemistry of biomolecules and regulation of biochemical pathways within and between cells, chemical biology remains distinct by focusing on the application of chemical tools to address biological questions.
History
Although considered a relatively new scientific field, the term "chemical biology" has been in use since the early 20th century, and has roots in scientific discovery from the early 19th century. The term 'chemical biology' can be traced back to an early appearance in a book published by Alonzo E. Taylor in 1907 titled "On Fermentation", and was subsequently used in John B. Leathes' 1930 article titled "The Harveian Oration on The Birth of Chemical Biology". However, it is unclear when the term was first used.
Friedrich Wöhler's 1828 synthesis of urea is an early example of the application of synthetic chemistry to advance biology. It showed that biological compounds could be synthesized with inorganic starting materials and weakened the previous notion of vitalism, or that a 'living' source was required to produce organic compounds. Wöhler's work is often considered to be instrumental in the development of organic chemistry and natural product synthesis, both of which play a large part in modern chemical biology.
Friedrich Miescher's work during the late 19th century investigating the cellular contents of human leukocytes led to the discovery of 'nuclein', which would later be renamed DNA. After isolating the nuclein from the nucleus of leukocytes through protease digestion, Miescher used chemical techniques such as elemental analysis and solubility tests to determine the composition of nuclein. This work would lay the foundations for Watson and Crick's discovery of the double-helix structure of DNA.
The rising interest in chemical biology has led to several journals dedicated to the field. Nature Chemical Biology, created in 2005, and ACS Chemical Biology, created in 2006, are two of the most well-known journals in this field, with impact factors of 14.8 and 4.0 respectively.
Nobel laureates in chemical biology
Research areas
Glycobiology
Glycobiology is the study of the structure and function of carbohydrates. While DNA, RNA, and proteins are encoded at the genetic level, carbohydrates are not encoded directly from the genome, and thus require different tools for their study. By applying chemical principles to glycobiology, novel methods for analyzing and synthesizing carbohydrates can be developed. For example, cells can be supplied with synthetic variants of natural sugars to probe their function. Carolyn Bertozzi's research group has developed methods for site-specifically reacting molecules at the surface of cells via synthetic sugars.
Combinatorial chemistry
Combinatorial chemistry involves simultaneously synthesizing a large number of related compounds for high-throughput analysis. Chemical biologists are able to use principles from combinatorial chemistry in synthesizing active drug compounds and maximizing screening efficiency. Similarly, these principles can be used in areas of agriculture and food research, specifically in the syntheses of unnatural products and in generating novel enzyme inhibitors.
Peptide synthesis
Chemical synthesis of proteins is a valuable tool in chemical biology as it allows for the introduction of non-natural amino acids as well as residue-specific incorporation of "posttranslational modifications" such as phosphorylation, glycosylation, acetylation, and even ubiquitination. These properties are valuable for chemical biologists as non-natural amino acids can be used to probe and alter the functionality of proteins, while post-translational modifications are widely known to regulate the structure and activity of proteins. Although strictly biological techniques have been developed to achieve these ends, the chemical synthesis of peptides often has a lower technical and practical barrier to obtaining small amounts of the desired protein.
To make protein-sized polypeptide chains with the small peptide fragments made by synthesis, chemical biologists can use the process of native chemical ligation. Native chemical ligation involves the coupling of a C-terminal thioester and an N-terminal cysteine residue, ultimately resulting in formation of a "native" amide bond. Other strategies that have been used for the ligation of peptide fragments using the acyl transfer chemistry first introduced with native chemical ligation include expressed protein ligation, sulfurization/desulfurization techniques, and use of removable thiol auxiliaries.
Enrichment techniques for proteomics
Chemical biologists work to improve proteomics through the development of enrichment strategies, chemical affinity tags, and new probes. Samples for proteomics often contain many peptide sequences and the sequence of interest may be highly represented or of low abundance, which creates a barrier for their detection. Chemical biology methods can reduce sample complexity by selective enrichment using affinity chromatography. This involves targeting a peptide with a distinguishing feature like a biotin label or a post translational modification. Methods have been developed that include the use of antibodies, lectins to capture glycoproteins, and immobilized metal ions to capture phosphorylated peptides and enzyme substrates to capture select enzymes.
Enzyme probes
To investigate enzymatic activity as opposed to total protein, activity-based reagents have been developed to label the enzymatically active form of proteins (see Activity-based proteomics). For example, serine hydrolase- and cysteine protease-inhibitors have been converted to suicide inhibitors. This strategy enhances the ability to selectively analyze low abundance constituents through direct targeting. Enzyme activity can also be monitored through converted substrate. Identification of enzyme substrates is a problem of significant difficulty in proteomics and is vital to the understanding of signal transduction pathways in cells. A method that has been developed uses "analog-sensitive" kinases to label substrates using an unnatural ATP analog, facilitating visualization and identification through a unique handle.
Employing biology
Many research programs are also focused on employing natural biomolecules to perform biological tasks or to support a new chemical method. In this regard, chemical biology researchers have shown that DNA can serve as a template for synthetic chemistry, self-assembling proteins can serve as a structural scaffold for new materials, and RNA can be evolved in vitro to produce new catalytic function. Additionally, heterobifunctional (two-sided) synthetic small molecules such as dimerizers or PROTACs bring two proteins together inside cells, which can synthetically induce important new biological functions such as targeted protein degradation.
Directed evolution
A primary goal of protein engineering is the design of novel peptides or proteins with a desired structure and chemical activity. Because our knowledge of the relationship between primary sequence, structure, and function of proteins is limited, rational design of new proteins with engineered activities is extremely challenging. In directed evolution, repeated cycles of genetic diversification followed by a screening or selection process, can be used to mimic natural selection in the laboratory to design new proteins with a desired activity.
Several methods exist for creating large libraries of sequence variants. Among the most widely used are subjecting DNA to UV radiation or chemical mutagens, error-prone PCR, degenerate codons, or recombination. Once a large library of variants is created, selection or screening techniques are used to find mutants with a desired attribute. Common selection/screening techniques include FACS, mRNA display, phage display, and in vitro compartmentalization. Once useful variants are found, their DNA sequence is amplified and subjected to further rounds of diversification and selection.
The development of directed evolution methods was honored in 2018 with the awarding of the Nobel Prize in Chemistry to Frances Arnold for evolution of enzymes, and George Smith and Gregory Winter for phage display.
Bioorthogonal reactions
Successful labeling of a molecule of interest requires specific functionalization of that molecule to react chemospecifically with an optical probe. For a labeling experiment to be considered robust, that functionalization must minimally perturb the system. Unfortunately, these requirements are often hard to meet. Many of the reactions normally available to organic chemists in the laboratory are unavailable in living systems. Water- and redox- sensitive reactions would not proceed, reagents prone to nucleophilic attack would offer no chemospecificity, and any reactions with large kinetic barriers would not find enough energy in the relatively low-heat environment of a living cell. Thus, chemists have recently developed a panel of bioorthogonal chemistry that proceed chemospecifically, despite the milieu of distracting reactive materials in vivo.
The coupling of a probe to a molecule of interest must occur within a reasonably short time frame; therefore, the kinetics of the coupling reaction should be highly favorable. Click chemistry is well suited to fill this niche, since click reactions are rapid, spontaneous, selective, and high-yielding. Unfortunately, the most famous "click reaction," a [3+2] cycloaddition between an azide and an acyclic alkyne, is copper-catalyzed, posing a serious problem for use in vivo due to copper's toxicity. To bypass the necessity for a catalyst, Carolyn R. Bertozzi's lab introduced inherent strain into the alkyne species by using a cyclic alkyne. In particular, cyclooctyne reacts with azido-molecules with distinctive vigor.
Discovery of biomolecules through metagenomics
The advances in modern sequencing technologies in the late 1990s allowed scientists to investigate DNA of communities of organisms in their natural environments ("eDNA"), without culturing individual species in the lab. This metagenomic approach enabled scientists to study a wide selection of organisms that were previously not characterized due in part to an incompetent growth condition. Sources of eDNA include soils, ocean, subsurface, hot springs, hydrothermal vents, polar ice caps, hypersaline habitats, and extreme pH environments. Of the many applications of metagenomics, researchers such as Jo Handelsman, Jon Clardy, and Robert M. Goodman, explored metagenomic approaches toward the discovery of biologically active molecules such as antibiotics.
Functional or homology screening strategies have been used to identify genes that produce small bioactive molecules. Functional metagenomic studies are designed to search for specific phenotypes that are associated with molecules with specific characteristics. Homology metagenomic studies, on the other hand, are designed to examine genes to identify conserved sequences that are previously associated with the expression of biologically active molecules.
Functional metagenomic studies enable the discovery of novel genes that encode biologically active molecules. These assays include top agar overlay assays where antibiotics generate zones of growth inhibition against test microbes, and pH assays that can screen for pH change due to newly synthesized molecules using pH indicator on an agar plate. Substrate-induced gene expression screening (SIGEX), a method to screen for the expression of genes that are induced by chemical compounds, has also been used to search for genes with specific functions. Homology-based metagenomic studies have led to a fast discovery of genes that have homologous sequences as the previously known genes that are responsible for the biosynthesis of biologically active molecules. As soon as the genes are sequenced, scientists can compare thousands of bacterial genomes simultaneously. The advantage over functional metagenomic assays is that homology metagenomic studies do not require a host organism system to express the metagenomes, thus this method can potentially save the time spent on analyzing nonfunctional genomes. These also led to the discovery of several novel proteins and small molecules. In addition, an in silico examination from the Global Ocean Metagenomic Survey found 20 new lantibiotic cyclases.
Kinases
Posttranslational modification of proteins with phosphate groups by kinases is a key regulatory step throughout all biological systems. Phosphorylation events, either phosphorylation by protein kinases or dephosphorylation by phosphatases, result in protein activation or deactivation. These events have an impact on the regulation of physiological pathways, which makes the ability to dissect and study these pathways integral to understanding the details of cellular processes. There exist a number of challenges—namely the sheer size of the phosphoproteome, the fleeting nature of phosphorylation events and related physical limitations of classical biological and biochemical techniques—that have limited the advancement of knowledge in this area.
Through the use of small molecule modulators of protein kinases, chemical biologists have gained a better understanding of the effects of protein phosphorylation. For example, nonselective and selective kinase inhibitors, such as a class of pyridinylimidazole compounds are potent inhibitors useful in the dissection of MAP kinase signaling pathways. These pyridinylimidazole compounds function by targeting the ATP binding pocket. Although this approach, as well as related approaches, with slight modifications, has proven effective in a number of cases, these compounds lack adequate specificity for more general applications. Another class of compounds, mechanism-based inhibitors, combines knowledge of the kinase enzymology with previously utilized inhibition motifs. For example, a "bisubstrate analog" inhibits kinase action by binding both the conserved ATP binding pocket and a protein/peptide recognition site on the specific kinase. Research groups also utilized ATP analogs as chemical probes to study kinases and identify their substrates.
The development of novel chemical means of incorporating phosphomimetic amino acids into proteins has provided important insight into the effects of phosphorylation events. Phosphorylation events have typically been studied by mutating an identified phosphorylation site (serine, threonine or tyrosine) to an amino acid, such as alanine, that cannot be phosphorylated. However, these techniques come with limitations and chemical biologists have developed improved ways of investigating protein phosphorylation. By installing phospho-serine, phospho-threonine or analogous phosphonate mimics into native proteins, researchers are able to perform in vivo studies to investigate the effects of phosphorylation by extending the amount of time a phosphorylation event occurs while minimizing the often-unfavorable effects of mutations. Expressed protein ligation, has proven to be successful techniques for synthetically producing proteins that contain phosphomimetic molecules at either terminus. In addition, researchers have used unnatural amino acid mutagenesis at targeted sites within a peptide sequence.
Advances in chemical biology have also improved upon classical techniques of imaging kinase action. For example, the development of peptide biosensors—peptides containing incorporated fluorophores improved temporal resolution of in vitro binding assays. One of the most useful techniques to study kinase action is Fluorescence Resonance Energy Transfer (FRET). To utilize FRET for phosphorylation studies, fluorescent proteins are coupled to both a phosphoamino acid binding domain and a peptide that can be phosphorylated. Upon phosphorylation or dephosphorylation of a substrate peptide, a conformational change occurs that results in a change in fluorescence. FRET has also been used in tandem with Fluorescence Lifetime Imaging Microscopy (FLIM) or fluorescently conjugated antibodies and flow cytometry to provide quantitative results with excellent temporal and spatial resolution.
Biological fluorescence
Chemical biologists often study the functions of biological macromolecules using fluorescence techniques. The advantage of fluorescence versus other techniques resides in its high sensitivity, non-invasiveness, safe detection, and ability to modulate the fluorescence signal. In recent years, the discovery of green fluorescent protein (GFP) by Roger Y. Tsien and others, hybrid systems and quantum dots have enabled assessing protein location and function more precisely. Three main types of fluorophores are used: small organic dyes, green fluorescent proteins, and quantum dots. Small organic dyes usually are less than 1 kDa, and have been modified to increase photostability and brightness, and reduce self-quenching. Quantum dots have very sharp wavelengths, high molar absorptivity and quantum yield. Both organic dyes and quantum dyes do not have the ability to recognize the protein of interest without the aid of antibodies, hence they must use immunolabeling. Fluorescent proteins are genetically encoded and can be fused to your protein of interest. Another genetic tagging technique is the tetracysteine biarsenical system, which requires modification of the targeted sequence that includes four cysteines, which binds membrane-permeable biarsenical molecules, the green and the red dyes "FlAsH" and "ReAsH", with picomolar affinity. Both fluorescent proteins and biarsenical tetracysteine can be expressed in live cells, but present major limitations in ectopic expression and might cause a loss of function.
Fluorescent techniques have been used to assess a number of protein dynamics including protein tracking, conformational changes, protein–protein interactions, protein synthesis and turnover, and enzyme activity, among others. Three general approaches for measuring protein net redistribution and diffusion are single-particle tracking, correlation spectroscopy and photomarking methods. In single-particle tracking, the individual molecule must be both bright and sparse enough to be tracked from one video to the other. Correlation spectroscopy analyzes the intensity fluctuations resulting from migration of fluorescent objects into and out of a small volume at the focus of a laser. In photomarking, a fluorescent protein can be dequenched in a subcellular area with the use of intense local illumination and the fate of the marked molecule can be imaged directly. Michalet and coworkers used quantum dots for single-particle tracking using biotin-quantum dots in HeLa cells. One of the best ways to detect conformational changes in proteins is to label the protein of interest with two fluorophores within close proximity. FRET will respond to internal conformational changes result from reorientation of one fluorophore with respect to the other. One can also use fluorescence to visualize enzyme activity, typically by using a quenched activity-based proteomics (qABP). Covalent binding of a qABP to the active site of the targeted enzyme will provide direct evidence concerning if the enzyme is responsible for the signal upon release of the quencher and regain of fluorescence.
Education in chemical biology
Undergraduate education
Despite an increase in biological research within chemistry departments, attempts at integrating chemical biology into undergraduate curricula are lacking. For example, although the American Chemical Society (ACS) requires for foundational courses in a Chemistry Bachelor's degree to include biochemistry, no other biology-related chemistry course is required.
Although a chemical biology course is often not required for an undergraduate degree in Chemistry, many universities now provide introductory chemical biology courses for their undergraduate students. The University of British Columbia, for example, offers a fourth-year course in synthetic chemical biology.
| Biology and health sciences | Biochemistry and molecular biology | null |
1686370 | https://en.wikipedia.org/wiki/Mineral%20processing | Mineral processing | Mineral processing is the process of separating commercially valuable minerals from their ores in the field of extractive metallurgy. Depending on the processes used in each instance, it is often referred to as ore dressing or ore milling.
Beneficiation is any process that improves (benefits) the economic value of the ore by removing the gangue minerals, which results in a higher grade product (ore concentrate) and a waste stream (tailings). There are many different types of beneficiation, with each step furthering the concentration of the original ore. Key is the concept of recovery, the mass (or equivalently molar) fraction of the valuable mineral (or metal) extracted from the ore and carried across to the concentrate.
History
Before the advent of heavy machinery, raw ore was broken up using hammers wielded by hand, a process called "spalling". Eventually, mechanical means were found to achieve this. For instance, stamp mills were being used in central Asia in the vicinity of Samarkand as early as 973. There is evidence the process was in use in Persia in the early medieval period. By the 11th century, stamp mills were in widespread use throughout the medieval Islamic world, from Islamic Spain and North Africa in the west to Central Asia in the east. A later example was the Cornish stamps, consisting of a series of iron hammers mounted in a vertical frame, raised by cams on the shaft of a waterwheel and falling onto the ore under gravity.
Iron beneficiation has been evident since as early as 800 BC in China with the use of bloomery. A bloomery is the original form of smelting and allowed people to make fires hot enough to melt oxides into a liquid that separates from the iron. Although the bloomery was promptly phased out by the invention of the blast furnace, it was still heavily relied on in Africa and Europe until the early part of the second millennium. The blast furnace was the next step in smelting iron which produced pig iron. The first blast furnaces in Europe appeared in the early 1200s around Sweden and Belgium, and not until the late 1400s in England. The pig iron poured from a blast furnace is high in carbon making it hard and brittle, making it hard to work with. In 1856 the Bessemer process was invented that turns the brittle pig iron into steel, a more malleable metal. Since then, many different technologies have been invented to replace the Bessemer process such as the electric arc furnace, basic oxygen steelmaking, and direct reduced iron (DRI).
For sulfide ores, a different process is taken for beneficiation. The ore needs to have the sulfur removed before smelting can begin. Roasting is the primary method of separating, where wood was placed on heaps of ore and set on fire to help with oxidation.
2 Cu2S + 3 O2 → 2 Cu2O + 2 SO2
The earliest practices of roasting were done outside, allowing large clouds of sulfur dioxide to blow over the land causing serious harm to surrounding ecosystems, both aquatic and terrestrial. The clouds of sulfur dioxide combined with local deforestation for wood needed for roasting compounded damages to the environment, as seen in Sudbury, Ontario and the Inco Superstack.
The simplest method of separating ore from the gangue consists of picking out the individual crystals of each. This is a very tedious process, particularly when the individual particles are small. Another comparatively simple method relies on the various minerals having different densities, causing them to collect in different places: metallic minerals (being heavier) will drop out of suspension more quickly than lighter ones, which will be carried further by a stream of water. The process of panning and sifting for gold uses both of these methods. Various devices known as 'bundles' were used to take advantage of this property. Later, more advanced machines were used such as the Frue vanner, invented in 1874.
Other equipment used historically includes the hutch, a trough used with some ore-dressing machines and the keeve or kieve, a large tub used for differential settlement.
Types of separation
Disaggregation
Beneficiation can begin within the mine itself. Most mines will have a crusher within the mine itself where separation of ore and gangue minerals occurs and as a side effect becomes easier to transport. After the crusher the ore will go through a grinder or a mill to get the ore into fine particles. Dense media separation (DMS) is used to further separate the desired ore from rocks and gangue minerals. This will stratify the crushed aggregate by density making separation easier. Where the DMS occurs in the process can be important, the grinders or mills will process much less waste rock if the DMS occurs beforehand. This will lower wear on the equipment as well as operating costs since there is a lower volume being put through.
Physical separation
After the milling stage the ore can be further separated from the rock. One way this can be achieved is by using the physical properties of the ore to separate it from the rest of the rock. Prior to any physical separation process, sizing of ore particles is important for effective separation. This is done by using either Industrial Screens or Classifiers.
These processes are gravity separation, flotation, and magnetic separation. Gravity separation uses centrifugal forces and specific gravity of ores and gangue to separate them. Magnetic separation is used to separate magnetic gangue from the desired ore, or conversely to remove a magnetic target ore from nonmagnetic gangue. DMS is also considered a physical separation.
Chemical separation
Some ore physical properties can not be relied on for separation, therefore chemical processes are used to separate the ores from the rock. Froth flotation, leaching, and electrowinning are the most common types of chemical separation. Froth flotation uses hydrophobic and hydrophilic properties to separate the ore from the gangue. Hydrophobic particles will rise to the top of the solution to be skimmed off. Changes to pH in the solution can influence what particles will be hydrophilic. Leaching works by dissolving the desired ore into solution from the rock. Electrowinning is not a primary method of separation, but is required to get the ore out of solution after leaching.
Unit operations
Mineral processing can involve four general types of unit operation:
1) Comminution – particle size reduction;
2) Sizing – separation of particle sizes by screening or classification;
3) Concentration by taking advantage of physical and surface chemical properties; and
4) Dewatering – solid/liquid separation.
In all of these processes, the most important considerations are the economics of the processes, which is dictated by the grade and recovery of the final product. To do this, the mineralogy of the ore needs to be considered as this dictates the amount of liberation required and the processes that can occur. The smaller the particles processes, the greater the theoretical grade and recovery of the final product. This, however, becomes difficult to do with fine particles since they prevent certain concentration processes from occurring.
Comminution
Comminution is particle size reduction of materials. Comminution may be carried out on either dry materials or slurries. Crushing and grinding are the two primary comminution processes. Crushing is normally carried out on run-of-mine ore, while grinding (normally carried out after crushing) may be conducted on dry or slurried material. In comminution, the size reduction of particles is done by three types of forces: compression, impact and attrition. Compression and impact forces are extensively used in crushing operations while attrition is the dominant force in grinding. The primarily used equipment in crushing are jaw crushers, gyratory crushers and cone crushers whereas rod mills and ball mills, usually closed circuited with a classifier unit, are generally employed for grinding purposes in a mineral processing plant. Crushing is a dry process whereas grinding is generally performed wet and hence is more energy intensive.
Sizing
Sizing is the general term for separation of particles according to their size. The simplest sizing process is screening, or passing the particles to be sized through a screen or number of screens. Screening equipment can include grizzlies, bar screens, wedge wire screens, radial sieves, banana screens, multi-deck screens, vibratory screen, fine screens, flip flop screens, and wire mesh screens. Screens can be static (typically the case for very coarse material), or they can incorporate mechanisms to shake or vibrate the screen. Some considerations in this process include the screen material, the aperture size, shape and orientation, the amount of near sized particles, the addition of water, the amplitude and frequency of the vibrations, the angle of inclination, the presence of harmful materials, like steel and wood, and the size distribution of the particles.
Classification refers to sizing operations that exploit the differences in settling velocities exhibited by particles of different size. Classification equipment may include ore sorters, gas cyclones, hydrocyclones, rotating trommels, rake classifiers or fluidized classifiers.
An important factor in both comminution and sizing operations is the determination of the particle size distribution of the materials being processed, commonly referred to as particle size analysis. Many techniques for analyzing particle size are used, and the techniques include both off-line analyses which require that a sample of the material be taken for analysis and on-line techniques that allow for analysis of the material as it flows through the process.
Concentration
There are a number of ways to increase the concentration of the wanted minerals: in any particular case, the method chosen will depend on the relative physical and surface chemical properties of the mineral and the gangue. Concentration is defined as the number of moles of a solute in a volume of the solution. In case of mineral processing, concentration means the increase of the percentage of the valuable mineral in the concentrate.
Gravity concentration
Gravity separation is the separation of two or more minerals of different specific gravity by their relative movement in response to the force of gravity and one or more other forces (such as centrifugal forces, magnetic forces, buoyant forces), one of which is resistance to motion (drag force) by a viscous medium such as heavy media, water or, less commonly, air.
Gravity separation is one of the oldest technique in mineral processing but has seen a decline in its use since the introduction of methods like flotation, classification, magnetic separation and leaching. Gravity separation dates back to at least 3000 BC when Egyptians used the technique for separation of gold.
It is necessary to determine the suitability of a gravity concentration process before it is employed for concentration of an ore. The concentration criterion is commonly used for this purpose, designated in the following equation (where represents specific gravity):
for CC > 2.5, suitable for separation of particles above 75 microns in size
for 1.75 < CC < 2.5, suitable for separation of particles above 150 microns in size
for 1.50 < CC < 1.75, suitable for separation of particles above 1.7 mm in size
for 1.25 < CC < 1.50, suitable for separation of particles above 6.35 mm in size
for CC < 1.25, not suitable for any size
Although concentration criteria is a useful rule of thumb when predicting amenability to gravity concentration, factors such as particle shape and relative concentration of heavy and light particles can dramatically affect separation efficiency in practice.
Classification
There are several methods that make use of the weight or density differences of particles:
Heavy media or dense media separation (these include baths, drums, larcodems, dyana whirlpool separators, and dense medium cyclones)
Shaking tables, such as the Wilfley table
Spiral separators
Reflux classifier
Jig concentrators are continuous processing gravity concentration devices using a pulsating fluidized bed (RMS-Ross Corp. circular jig plants)
Centrifugal bowl concentrators, such as the Knelson concentrator
Multi-gravity separators including Knelson, Mozley (multi or enhanced) gravity separator, Salter cyclones (multi-gravity separator) and the Kelsey jig)
Inline pressure jigs
Reichert cones
Sluices
Elutriators
Optima classifier- a water -only density separator. It is an engineered system that consists of feed preparation, beneficiation/separation and dewatering modules to ensure a high-quality product with the correct moisture content. Typical products consist of ash below 16% and a calorific value of at least 27-million joules per kilogram.
These processes can be classified as either density separation or gravity (weight) separation. In dense media separation a media is created with a density in between the density of the ore and gangue particles. When subjected to this media particles either float or sink depending on their density relative to the media. In this way the separation takes place purely on density differences and does not, in principle, relay on any other factors such as particle weight or shape. In practice, particle size and shape can affect separation efficiency. Dense medium separation can be performed using a variety of mediums. These include, organic liquids, aqueous solutions or suspensions of very fine particles in water or air. The organic liquids are typically not used due to their toxicity, difficulties in handling and relative cost. Industrially, the most common dense media is a suspension of fine magnetite and/or ferrosilicon particles. An aqueous solution as a dense medium is used in coal processing in the form of a belknap wash and suspensions in air are used in water-deficient areas, like areas of China, where sand is used to separate coal from the gangue minerals.
Gravity separation is also called relative gravity separation as it separates particles due to their relative response to a driving force. This is controlled by factors such as particle weight, size and shape. These processes can also be classified into multi-G and single G processes. The difference is the magnitude of the driving force for the separation. Multi-G processes allow the separation of very fine particles to occur (in the range of 5 to 50 micron) by increasing the driving force of separation in order to increase the rate at which particles separate. In general, single G process are only capable of processing particles that are greater than approximately 80 micron in diameter.
Of the gravity separation processes, the spiral concentrators and circular jigs are two of the most economical due to their simplicity and use of space. They operate by flowing film separation and can either use washwater or be washwater-less. The washwater spirals separate particles more easily but can have issues with entrainment of gangue with the concentrate produced.
Froth flotation
Froth flotation is an important concentration process. This process can be used to separate any two different particles and operated by the surface chemistry of the particles. In flotation, bubbles are introduced into a pulp and the bubbles rise through the pulp. In the process, hydrophobic particles become bound to the surface of the bubbles. The driving force for this attachment is the change in the surface free energy when the attachment occurs. These bubbles rise through the slurry and are collected from the surface. To enable these particles to attach, careful consideration of the chemistry of the pulp needs to be made. These considerations include the pH, Eh and the presence of flotation reagents. The pH is important as it changes the charge of the particles surface and the pH affects the chemisorption of collectors on the surface of the particles.
The addition of flotation reagents also affects the operation of these processes. The most important chemical that is added is the collector. This chemical binds to the surface of the particles as it is a surfactant. The main considerations in this chemical is the nature of the head group and the size of the hydrocarbon chain. The hydrocarbon tail needs to be short to maximize the selectivity of the desired mineral and the headgroup dictates which minerals it attaches to.
The frothers are another important chemical addition to the pulp or slurry as they enable stable bubbles to be formed. This is important because if the bubbles coalesce, minerals will fall off their surface. The bubbles however should not be too stable as this prevents easy transportation and dewatering of the concentrate formed. The mechanism of these frothers is not completely known and further research into their mechanisms is being performed.
Depressants and activators are used to selectively separate one mineral from another. Depressants inhibit the flotation of one mineral or minerals while activators enable the flotation of others. Examples of these include CN−, used to depress all sulfides but galena and this depressant is believed to operate by changing the solubility of chemisorbed and physisorbed collectors on sulfides. This theory originates from Russia. An example of an activator is Cu2+ ions, used for the flotation of sphalerite.
There are a number of cells able to be used for the flotation of minerals. these include flotation columns and mechanical flotation cells. The flotation columns are used for finer minerals and typically have a higher grade and lower recovery of minerals than mechanical flotation cells. The cells in use at the moment can exceed 300 m3. This is done as they are cheaper per unit volume than smaller cells, but they are not able to be controlled as easily as smaller cells.
This process was invented in the 19th century in Australia. It was used to recover a sphalerite concentrate from tailings, produced using gravity concentration. Further improvements have come from Australia in the form of the Jameson Cell, developed at the University of Newcastle, Australia. This operated by the use of a plunging jet that generates fine bubbles. These fine bubbles have a higher kinetic energy and as such they can be used for the flotation of fine grained minerals, such as those produced by the IsaMill.
Staged flotation reactors (SFRs) split the flotation process into three defined stages per cell. They are becoming increasingly more common in use as they require much less energy, air and installation space.
Electrostatic separation
There are two main types of electrostatic separators. These work in similar ways, but the forces applied to the particles are different and these forces are gravity and electrostatic attraction. The two types are electrodynamic separators (or high tension rollers) or electrostatic separators. In high tension rollers, particles are charged by a corona discharge. This charges the particles that subsequently travel on a drum. The conducting particles lose their charge to the drum and are removed from the drum with centripetal acceleration. Electrostatic plate separators work by passing a stream of particles past a charged anode. The conductors lose electrons to the plate and are pulled away from the other particles due to the induced attraction to the anode. These separators are used for particles between 75 and 250 micron and for efficient separation to occur, the particles need to be dry, have a close size distribution and uniform in shape. Of these considerations, one of the most important is the water content of the particles. This is important as a layer of moisture on the particles will render the non-conductors as conductors as the layer of the water is conductive.
Electrostatic plate separators are usually used for streams that have small conductors and coarse non-conductors. The high tension rollers are usually used for streams that have coarse conductors and fine non-conductors.
These separators are commonly used for separating mineral sands, an example of one of these mineral processing plants is the CRL processing plant at Pinkenba in Brisbane Queensland. In this plant, zircon, rutile and ilmenite are separated from the silica gangue. In this plant, the separation is performed in a number of stages with roughers, cleaners, scavengers and recleaners.
Magnetic separation
Magnetic separation is a process in which magnetically susceptible material is extracted from a mixture using a magnetic force. This separation technique can be useful in mining iron as it is attracted to a magnet. In mines where wolframite was mixed with cassiterite, such as South Crofty and East Pool mine in Cornwall or with bismuth such as at the Shepherd and Murphy mine in Moina, Tasmania, magnetic separation was used to separate the ores. At these mines a device called a Wetherill's Magnetic Separator (invented by John Price Wetherill, 1844–1906)[1] was used. In this machine the raw ore, after calcination was fed onto a moving belt which passed underneath two pairs of electromagnets under which further belts ran at right angles to the feed belt. The first pair of electromagnets was weakly magnetised and served to draw off any iron ore present. The second pair were strongly magnetised and attracted the wolframite, which is weakly magnetic. These machines were capable of treating 10 tons of ore a day. This process of separating magnetic substances from the non-magnetic substances in a mixture with the help of a magnet is called magnetic separation..
This process operates by moving particles in a magnetic field. The force experienced in the magnetic field is given by the equation f=m/k.H.dh/dx. with k=magnetic susceptibility, H-magnetic field strength, and dh/dx being the magnetic field gradient. As seen in this equation, the separation can be driven in two ways, either through a gradient in a magnetic field or the strength of a magnetic field. The different driving forces are used in the different concentrators. These can be either with water or without. Like the spirals, washwater aids in the separation of the particles while increases the entrainment of the gangue in the concentrate.
Automated Ore Sorting
Modern, automated sorting applies optical sensors (visible spectrum, near infrared, X-ray, ultraviolet), that can be coupled with electrical conductivity and magnetic susceptibility sensors, to control the mechanical separation of ore into two or more categories on an individual rock by rock basis. Also new sensors have been developed which exploit material properties such as electrical conductivity, magnetization, molecular structure and thermal conductivity. Sensor based sorting has found application in the processing of nickel, gold, copper, coal
and diamonds.
Dewatering
Dewatering is an important process in mineral processing. The purpose of dewatering is to remove water absorbed by the particles which increases the pulp density. This is done for a number of reasons, specifically, to enable ore handling and concentrates to be transported easily, allow further processing to occur and to dispose of the gangue. The water extracted from the ore by dewatering is recirculated for plant operations after being sent to a water treatment plant. The main processes that are used in dewatering include dewatering screens, sedimentation, filtering, and thermal drying. These processes increase in difficulty and cost as the particle size decreases.
Dewatering screens operate by passing particles over a screen. The particles pass over the screen while the water passes through the apertures in the screen. This process is only viable for coarse ores that have a close size distribution as the apertures can allow small particles to pass through.
Sedimentation operates by passing water into a large thickener or clarifier. In these devices, the particles settle out of the slurry under the effects of gravity, or centripetal forces. These are limited by the surface chemistry of the particles and the size of the particles. To aid in the sedimentation process, flocculants and coagulants are added to reduce the repulsive forces between the particles. This repulsive force is due to the double layer formed on the surface of the particles. The flocculants work by binding multiple particles together while the coagulants work by reducing the thickness of the charged layer on the outside of the particle. After thickening, slurry is often stored in ponds or impoundments. Alternatively, it can pumped into a belt press or membrane filter press to recycle process water and create stackable, dry filter cake, or "tailings".
Thermal drying is usually used for fine particles and to remove low water content in the particles. Some common processes include rotary dryers, fluidized beds, spray driers, hearth dryers and rotary tray dryers. This process is usually expensive to operate due to the fuel requirement of the dryers.
Other processes
Many mechanical plants also incorporate hydrometallurgical or pyrometallurgical processes as part of an extractive metallurgical operation. Geometallurgy is a branch of extractive metallurgy that combines mineral processing with the geologic sciences. This includes the study of oil agglomeration
A number of auxiliary materials handling operations are also considered a branch of mineral processing such as storage (as in bin design), conveying, sampling, weighing, slurry transport, and pneumatic transport.
The efficiency and efficacy of many processing techniques are influenced by upstream activities such as mining method and blending.
Case examples
In the case of gold, after adsorbing onto carbon, it is put into a sodium hydroxide and cyanide solution. In the solution the gold is pulled out of the carbon and into the solution. The gold ions are removed from solution at steel wool cathodes from electrowinning. The gold then goes off to be smelted.
Lithium is hard to separate from gangue due to similarities in the minerals. In order to separate the lithium both physical and chemical separation techniques are used. First froth flotation is used. Due to similarities in mineralogy there is not complete separation after flotation. The gangue that is found with lithium after the flotation are often iron bearing. The float concentrate goes through magnetic separation to remove the magnetic gangue from the nonmagnetic lithium.
Conferences
European Metallurgical Conference (EMC)
EMC, the European Metallurgical Conference has developed to the most important networking business event dedicated to the non-ferrous metals industry in Europe. From the start of the conference sequence in 2001 at Friedrichshafen it was host to some of most relevant metallurgists from all countries of the world. The conference is held every two years by invitation of GDMB Society of Metallurgists and Miners and is particularly directed to metal producers, plant manufactures, equipment suppliers and service providers as well as members of universities and consultants.
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1686418 | https://en.wikipedia.org/wiki/Voltage%20source | Voltage source | A voltage source is a two-terminal device which can maintain a fixed voltage. An ideal voltage source can maintain the fixed voltage independent of the load resistance or the output current. However, a real-world voltage source cannot supply unlimited current.
A voltage source is the dual of a current source. Real-world sources of electrical energy, such as batteries and generators, can be modeled for analysis purposes as a combination of an ideal voltage source and additional combinations of impedance elements.
Ideal voltage sources
An ideal voltage source is a two-terminal device that maintains a fixed voltage drop across its terminals. It is often used as a mathematical abstraction that simplifies the analysis of real electric circuits. If the voltage across an ideal voltage source can be specified independently of any other variable in a circuit, it is called an independent voltage source. Conversely, if the voltage across an ideal voltage source is determined by some other voltage or current in a circuit, it is called a dependent or controlled voltage source. A mathematical model of an amplifier will include dependent voltage sources whose magnitude is governed by some fixed relation to an input signal, for example. In the analysis of faults on electrical power systems, the whole network of interconnected sources and transmission lines can be usefully replaced by an ideal (AC) voltage source and a single equivalent impedance.
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The internal resistance of an ideal voltage source is zero; it is able to supply or absorb any amount of current. The current through an ideal voltage source is completely determined by the external circuit. When connected to an open circuit, there is zero current and thus zero power. When connected to a load resistance, the current through the source approaches infinity as the load resistance approaches zero (a short circuit). Thus, an ideal voltage source can supply unlimited power.
If two ideal independent voltage source are directly connected in parallel, they must have exactly the same voltage; Otherwise, it creates a fallacy in logic, similar to writing down the equation .
If two ideal independent voltage sources are connected in parallel via a resistor, the source with the lower voltage becomes a consumer. If the source voltages are equal, there is zero current and thus zero power.
Voltage sources in parallel shares the burden of current: If an exact duplicate of voltage is connected in parallel to the original one, either one of them will provide half of the electric current that the original voltage source would provide. For the remainder of the circuit, nothing has changed: These two voltage sources together provide the same voltage, and the same current as the original one alone.
No real voltage source is ideal; all have a non-zero effective internal resistance, and none can supply unlimited current. However, the internal resistance of a real voltage source is effectively modeled in linear circuit analysis by combining a non-zero resistance in series with an ideal voltage source (a Thévenin equivalent circuit).
Comparison between voltage and current sources
Most sources of electrical energy (the mains, a battery) are modeled as voltage sources.
An ideal voltage source provides no energy when it is loaded by an open circuit (i.e. an infinite impedance), but approaches infinite energy and current when the load resistance approaches zero (a short circuit). Such a theoretical device would have a zero ohm output impedance in series with the source. A real-world voltage source has a very low, but non-zero internal resistance and output impedance, often much less than 1 ohm.
Conversely, a current source provides a constant current, as long as the load connected to the source terminals has sufficiently low impedance. An ideal current source would provide no energy to a short circuit and approach infinite energy and voltage as the load resistance approaches infinity (an open circuit). An ideal current source has an infinite output impedance in parallel with the source. A real-world current source has a very high, but finite output impedance. In the case of transistor current sources, impedance of a few megohms (at low frequencies) is typical.
Since no ideal sources of either variety exist (all real-world examples have finite and non-zero source impedance), any current source can be considered as a voltage source with the same source impedance and vice versa. Voltage sources and current sources are sometimes said to be duals of each other and any non ideal source can be converted from one to the other by applying Norton's theorem or Thévenin's theorem.
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1687186 | https://en.wikipedia.org/wiki/Inland%20taipan | Inland taipan | The inland taipan (Oxyuranus microlepidotus), also commonly known as the western taipan, small-scaled snake, or fierce snake, is a species of extremely venomous snake in the family Elapidae. The species is endemic to semiarid regions of central east Australia. Aboriginal Australians living in those regions named the snake dandarabilla. It was formally described by Frederick McCoy in 1879 and then by William John Macleay in 1882, but for the next 90 years, it was a mystery to the scientific community; no further specimens were found, and virtually nothing was added to the knowledge of this species until its rediscovery in 1972.
Based on the median lethal dose value in mice, the venom of the inland taipan is by far the most toxic of any snake – much more so than even that of sea snakes – and it has the most toxic venom of any reptile when tested on human heart cell culture. The inland taipan is a specialist hunter of mammals, so its venom is specially adapted to kill warm-blooded species. One bite possesses enough lethality to kill more than an estimated 100 fully grown humans. It is an extremely fast and agile snake that can strike instantly with extreme accuracy, often striking multiple times in the same attack, and it envenomates in almost every case.
Although the most venomous and a capable striker, in contrast to the coastal taipan, which many experts cite as an extremely dangerous snake due to its behaviour when it encounters humans, the inland taipan is usually quite a shy and reclusive snake, with a placid disposition, and prefers to escape from trouble. However, it will defend itself and strike if provoked, mishandled, or prevented from escaping. Because it lives in such remote locations, the inland taipan seldom comes in contact with people; therefore it is not considered the deadliest snake in the world overall, especially in terms of disposition and human deaths per year. The word "fierce" from its alternative name describes its venom, not its temperament.
Taxonomy
To the Aboriginal people from the place now called Goyder Lagoon in north-east South Australia, the inland taipan was called dandarabilla.
The inland taipan was first described scientifically in 1879. Two specimens of the fierce snake were discovered at the junction of the Murray and Darling Rivers in northwestern Victoria and described by Frederick McCoy, who called the species Diemenia microlepidota, or small-scaled brown snake. In 1882, a third specimen was found near Bourke, New South Wales, and William John Macleay described the same snake under the name Diemenia ferox (thinking it was a different species). No more specimens were collected until 1972. In 1896, George Albert Boulenger classified both as belonging to the same genus, Pseudechis (black snakes), referring to them as Pseudechis microlepidotus and P. ferox.
In 1956, relying only on published descriptions and notes, James Roy Kinghorn regarded ferox as a synonym for microlepidotus and proposed the genus Parademansia. In 1963, Eric Worrell considered Parademansia microlepidotus and Oxyuranus scutellatus (coastal taipan, named simply "taipan" in those days) to be the same species.
In September 1972, after receiving an unclassified snake head sample from a grazier from one of the Channel Country stations west of Windorah of the far southwest Queensland, herpetologists Jeanette Covacevich (then working for the Queensland Museum) and Charles Tanner travelled to the site and found 13 living specimens, and rediscovered the lost snake Parademansia microlepidotus. In 1976, Covacevich and Wombey argued that Parademansia microlepidotus belongs to a distinct genus, and this was also the opinion of Harold Cogger.
Covacevich, McDowell, Tanner & Mengden (1981) successfully argued, by comparing anatomical features, chromosomes, and behaviours of the two species then known as Oxyuranus scutellatus (taipan) and Parademansia microlepidota, that they belonged in a single genus. Oxyuranus (1923), the more senior name, was adopted for the combined genus.
Oxyuranus microlepidotus has been the fierce snake's binomial name since the early 1980s. The generic name Oxyuranus is from Greek oxys "sharp, needle-like", and ouranos "an arch" (specifically the arch of the heavens) and refers to the needle-like anterior process on the arch of the palate. The specific name microlepidotus means "small-scaled" (Latin). Hence the common name, "small-scaled snake". Since Covacevich et al., 1981, determined that the fierce snake (formerly: Parademansia microlepidota) is actually part of the genus Oxyuranus (taipan), another species, Oxyuranus scutellatus, which was previously known simply as the "taipan" (coined from the aboriginal snake's name dhayban), was renamed the "coastal taipan" (or "eastern taipan"), while the now newly classified Oxyuranus microlepidotus became commonly known as the "inland taipan" (or "western taipan").
Distribution and habitat
The inland taipan inhabits the black soil plains in the semiarid regions where the Queensland and South Australia borders converge.
In Queensland, the snake has been observed in Channel Country region (e.g., Diamantina National Park, Durrie Station, Morney Plains Station and Astrebla Downs National Park) and in South Australia it has been observed in the Marree-Innamincka NRM District (e.g., Goyder Lagoon Tirari Desert, Sturt Stony Desert, Coongie Lakes, Innamincka Regional Reserve and Oodnadatta). An isolated population also occurs near Coober Pedy, South Australia.
Two old records exist for localities further south-east, i.e., the junction of the Murray and Darling Rivers in northwestern Victoria (1879) and Bourke, New South Wales (1882), but the species has not been observed in either state since then.
Conservation status
Like every Australian snake, the inland taipan is protected by law.
Conservation status for the snake was assessed for the IUCN Red List for the first time in July 2017, and in 2018, was designated as least concern, stating, "This species is listed as least concern, as it is widespread and overall, it is not considered to be declining. Although the impact of potential threats requires further research, these are likely to be localized within the snake's range."
The inland taipan's conservation status has also been designated by Australian official sources:
South Australia: (Outback regional status) - least concern
Queensland: Rare (before 2010), near threatened (May 2010-December 2014), least concern (December 2014 – present)
New South Wales: Presumed extinct, because it "hasn't been recorded in its habitat...despite surveys in a time frame appropriate to their life cycle and type"
Victoria: Regionally extinct, based on the criterion: "As for Extinct but within a defined region (in this case the state of Victoria) that does not encompass the entire geographic range of the taxon. A taxon is presumed Regionally Extinct when exhaustive surveys in known and/or expected habitat, at appropriate ), throughout the region have failed to record an individual. Surveys should be over a time frame appropriate to the taxon’s life cycle and life form." The Australian Museum lists it as presumed extinct.
In captivity
Inland taipans are held in several zoo collections in Australia and overseas including the Adelaide Zoo and Taronga Zoo in Sydney.
The inland taipan is also on public display in Australia at the Australia Zoo, Australian Reptile Park, Billabong Sanctuary, Cairns Tropical Zoo, Lone Pine Koala Sanctuary and Shoalhaven Zoo.
The snake is also on display at several locations outside of Australia. In the United States, inland taipans are held at the Reptile Gardens in South Dakota, at Kentucky Reptile Zoo and at Animal World & Snake Farm Zoo in Texas.
In Europe, inland taipans are held in Sweden at the Stockholm Skansen Zoo and Gothenburg Universeum Moscow Zoo in Russia, (in the Moscow Zoo they are kept in the House of Reptiles, which is not usually open to the general public.) and in the UK at the London Zoo. Amateur zoo listings also report the snake at the tropicarium park in Jesolo, Italy, in Gifttierhaus Eimsheim, Welt der Gifte Greifswald, and Terra Zoo Rheinberg, Germany, in Lausanne vivarium Lausanne, Switzerland, in Randers Tropical Zoo Denmark, in Plzeň Zoo Czech Republic, and in Reptilienzoo Nockalm Patergassen Austria.
In Asia, inland taipans are held in the Singapore Zoo.
Private ownership law
In New South Wales, private ownership of an inland taipan is legal only with the highest class of venomous reptile licence.
Description
The inland taipan is dark tan, ranging from a rich, dark hue to a brownish light green, depending on the season. Its back, sides, and tail may be different shades of brown and grey, with many scales having a wide, blackish edge. These dark-marked scales occur in diagonal rows so that the marks align to form broken chevrons of variable length that are inclined backward and downward. The lowermost lateral scales often have an anterior yellow edge. The dorsal scales are smooth and without keels. The round-snouted head and neck are usually noticeably darker than the body (glossy black in winter; dark brown in summer), the darker colour allowing the snake to heat itself while exposing only a smaller portion of the body at the burrow entrance. The eyes are of average size with blackish-brown irises and without a noticeable coloured rim around the pupils.
It has 23 rows of dorsal scales at midbody, between 55 and 70 divided subcaudal scales, and one anal scale.
The inland taipan averages about in total length, although larger specimens can reach total lengths of . Its fangs are between 3.5 and 6.2 mm long (shorter than those of the coastal taipan).
Seasonal adaptation
Inland taipans adapt to their environments by changing the colour of their skin during seasonal changes. They tend to become lighter during the summer and darker during the winter. This seasonal colour change facilitates thermoregulation, allowing the snake to absorb more radiant heat in the colder months.
Breeding
Inland taipans produce clutches of one to two dozen eggs. The eggs hatch in about two months. They are usually laid in abandoned animal burrows and deep crevices. Reproduction rate depend in part on their diet: if not enough food is available, then the snake reproduces less.
Captive snakes generally live for 10 to 15 years. An inland taipan at the Australia Zoo lived to be over 20 years old.
Feeding
In the wild, the inland taipan consumes only mammals, mostly rodents, such as the long-haired rat (Rattus villosissimus), the plains rat (Pseudomys australis), the introduced house mouse (Mus musculus), and other dasyurids. In captivity, it may also eat day-old chicks. Unlike other venomous snakes that strike with a single, accurate bite then retreat while waiting for the prey to die, the fierce snake subdues the prey with a series of rapid, accurate strikes. It is known to deliver up to eight venomous bites in a single attack, often snapping its jaws fiercely several times to inflict multiple punctures in the same attack. Its more risky attack strategy entails holding its prey with its body and biting it repeatedly. This injects the extremely toxic venom deep into the prey. The venom acts so rapidly that its prey does not have time to fight back.
Natural threats
The mulga snake (Pseudechis australis) is immune to most Australian snake venom, and is known to also eat young inland taipans. The perentie (Varanus giganteus), a large monitor lizard, shares the same habitat. As it grows large enough, it readily tackles large venomous snakes as prey.
Interaction with humans
Many reptile keepers consider it a placid snake to handle.
Inland taipans are rarely encountered in the wild by the average person because of their remoteness and brief above-ground appearance during the day. So long as a person is not creating much vibration and noise, the inland taipan may not feel alarmed or bothered by a human presence. However, caution should be exercised and a safe distance maintained as it can inflict a potentially fatal bite. The inland taipan will defend itself and strike if provoked, mishandled, or prevented from escaping. Firstly, but not always, it makes a threat display by raising its forebody in a tight low S-shaped curve with its head facing the threat. Should the person choose to ignore the warning, the inland taipan will strike. It is an extremely fast and agile snake that can strike instantly with extreme accuracy, and it envenoms in almost every case.
Clinical toxicologist, venom researcher, herpetologist, and family physician Scott A. Weinstein et al. have stated in Toxicon journal (October 2017) "There have been 11 previously well-documented envenomings by O. microlepidotus, but only 2 were inflicted by wild snakes. When clinically indicated, prompt provision of adequate antivenom is the cornerstone of managing O. microlepidotus envenoming. Rapid application of pressure-bandage immobilization and efficient retrieval of victims envenomed in remote locales, preferably by medically well-equipped aircraft, probably improves the likelihood of a positive outcome."
Snakebite victims
A case of survival without antivenom was recorded in 1967; on 15 September, a tour guide was bitten while trying to capture a snake for a tour group in the Channel Country. He was conveyed to Broken Hill Hospital and then to Queen Elizabeth Hospital in Adelaide, but was not given antivenom, as he reported he was severely allergic to horse serum and believed he had been bitten by a brown snake. He spent four weeks in the hospital overall; his condition was likened to severe myasthenia gravis. Meanwhile, the snake was sent to Eric Worrell, who confirmed it was a coastal taipan. After its rediscovery in 1972, it was identified as an inland taipan.
In 1984, Australian toxicologist Peter Mirtschin was bitten by a 3-week-old inland taipan. He was the first to be treated with Taipan antivenom.
In September 2012, in the small city of Kurri Kurri, New South Wales, north of Sydney, more than 1000 km away from the snake's natural environment, a teenaged boy was bitten on the finger by an inland taipan. The teenager's rapid self-application of a compression bandage above the wound and the availability and administration of a polyvalent (broad-spectrum) antivenom in the local hospital saved his life. The police worked to find out how the inland taipan got to this part of Australia. The snake was most likely a stolen or illegal pet and the boy had tried to feed it.
In December 2013, reptile handler Scott Grant (age 40+), who was conducting a demonstration in front of 300 people at the annual building union's picnic in Portland, Victoria, had just finished showing the crowd an inland taipan and was trying to put it into a bag when it struck him. He got into his utility and tied a bandage around his arm. A few minutes later, however, he was lying on the ground and convulsing. He was flown in a serious condition to Essendon Airport and driven to the Royal Melbourne Hospital, where his condition was stabilised, and over time, he recovered. Only a tiny amount of venom from the inland taipan had entered his body, and the adverse reaction he felt shortly after was an allergic one, presumably due to his past snake bites.
In October 2017, Weinstein et al. published a case report in Toxicon, writing, "The victim was seeking to observe members of an isolated population of this species and was envenomed while attempting to photograph an approximately 1.5 m specimen. He reported feeling “drowsiness” and blurred vision that progressed to ptosis; he later developed dysphagia and dysarthria. The patient was treated with one vial of polyvalent antivenom, which was later followed by an additional two vials of taipan monovalent. He was intubated during retrieval, and recovered after 3 days of intensive care. He had a right ophthalmoplegia that persisted for approximately 1 week post-envenoming.".
According to Rob Bredl, "the Barefoot Bushman", in an isolated area of South Australia, his father, Joe Bredl, was bitten while catching an inland taipan and barely survived. A more recent victim was his friend John Robinson, bitten while cleaning an inland taipan's cage at his reptile display on the Sunshine Coast, Queensland. He weathered the bite without antivenom, but sustained considerable muscle damage, as well as heart damage.
Almost all positively identified inland taipan bite victims have been herpetologists handling the snakes for study or snake handlers, such as people who catch snakes to extract their venom, or keepers in wildlife parks. All were treated successfully with antivenom. No recorded incidents have been fatal since the advent of monovalent (specific) antivenom therapy, though weeks are needed to recover from such a severe bite.
Jeff Leibowitz, an American snake collector, was bitten by an inland taipan in 2024. Leibowitz was hospitalized for almost a week in the hospital in critical condition, given anti-venom, and survived.
Venom
The average quantity of venom delivered by this species is 44 mg, and the maximum dose recorded is 110 mg, compared to the Indian cobra (Naja naja) 169 mg/max 610 mg, and the North American eastern diamondback rattlesnake (Crotalus adamanteus) 410 mg/max 848 mg.
The median lethal dose (LD50), subcutaneous (the most applicable to actual bites) for mice is 0.025 mg/kg (0.01 mg/kg subcutaneous, in bovine serum albumin). Compared to the beaked sea snake (Enhydrina schistosa) 0.164 mg/kg, Indian cobra 0.565 mg/kg, North American eastern diamondback rattlesnake 11.4 mg/kg, the inland taipan has a smaller venom yield than its cousin the coastal taipan yet its venom is almost four times as toxic. One bite's worth of venom is enough to kill 100 fully grown men.
Intravenous, intraperitoneal, and intramuscular LD50s for the inland taipan venom have not been tested.
Belcher's sea snake (Hydrophis belcheri), which many times is mistakenly called the hook-nosed sea snake (Enhydrina schistosa), has been erroneously popularized as the most venomous snake in the world, due to Ernst and Zug's published book Snakes in Question: The Smithsonian Answer Book from 1996. Bryan Grieg Fry, a prominent venom expert, has clarified the error: "The hook-nosed myth was due to a fundamental error in a book called Snakes in Question. In there, all the toxicity testing results were lumped in together, regardless of the mode of testing (e.g. subcutaneous vs. intramuscular vs. intravenous vs. intraperitoneal). As the mode can influence the relative number, venoms can only be compared within a mode. Otherwise, it's apples and rocks." Belcher's sea snake's actual LD50 (recorded only intramuscularly) is 0.24 mg/kg and 0.155 mg/kg, less lethal than other sea snakes such as the olive sea snake (Aipysurus laevis) 0.09 mg/kg and the most toxic intramuscularly, recorded of the sea snakes – the black-banded robust sea snake (Hydrophis melanosoma) 0.082 mg/kg. The black-banded robust sea snake has also been tested subcutaneously registering at 0.111 mg/kg, which is in line with the coastal taipan and thus more than four times less toxic than the inland taipan's venom. In the LD50 subcutaneous test, it is actually Dubois' sea snake (Aipysurus duboisii) which has the most toxic venom of any of the sea snakes tested, registering at 0.044 mg/kg. This is still nearly half as lethal as the inland taipan's venom.
The biological properties and toxicity of a baby inland taipan's venom are not significantly different from or weaker than those of an adult's.
The inland taipan's venom consists of:
Neurotoxins: Presynaptic neurotoxins; paradoxin (PDX), and postsynaptic neurotoxins; Oxylepitoxin-1, alpha-oxytoxin 1, alpha-scutoxin 1 – affecting the nervous system.
Hemotoxins (procoagulants) – affecting the blood
Myotoxins – affecting the muscles
Possibly nephrotoxins – affecting the kidneys
Possibly haemorrhagins – affecting the blood vessels (endothelium)
Hyaluronidase enzyme – increases the rate of absorption of venom
Paradoxin (PDX) appears to be one of the most potent, if not the most potent, beta-neurotoxins yet discovered. Beta-neurotoxins keep nerve endings from liberating the neurotransmitter acetylcholine.
According to researcher Ronelle Welton of James Cook University, most of the contents in the venom have not been characterized and little molecular research has been undertaken on taipan (Oxyuranus) species at large. As of 2005, the amino acid sequences of only seven proteins from inland taipan have been submitted to SWISS-PROT databases.
Clinical effects
The mortality rate is high in untreated cases:
Dangerousness of bite - severe envenomation likely, high lethality potential
Rate of envenoming - >80%
Untreated lethality rate - >80%
Clinically, envenomation may represent a complex scenario of multiple organ-system poisoning, with neurotoxic symptoms typically dominating. Acute kidney injury, rhabdomyolysis, and disseminated coagulopathy may also complicate the setting.
The first local and general symptoms of a bite are local pain and variable nonspecific effects, which may include headache, nausea, vomiting, abdominal pain, diarrhoea, dizziness, collapse, or convulsions leading to major organ effects - neurotoxicity, coagulopathy, rhabdomyolysis or kidney failure/damage, and finally death.
Inland taipan snake venom contains potent presynaptic neurotoxins (toxins in venom that cause paralysis or muscle weakness). Also present are postsynaptic neurotoxins, which are less potent, but more rapidly acting than the presynaptic neurotoxins. Presynaptic neurotoxins disrupt neurotransmitter release from the axon terminal. This takes days to resolve and does not respond to antivenom. Postsynaptic neurotoxins competitively block acetylcholine receptors, but the effect can be reversed by antivenom. Envenoming causes a progressive descending flaccid paralysis; ptosis is usually the first sign, then facial (dysarthria) and bulbar involvement occur, progressing to dyspnea and respiratory paralysis leading to suffocation and peripheral weakness. Because it can act so fast, it can kill a person within about 45 minutes. People experiencing effects of the venom within half an hour have been reported. The development of general or respiratory paralysis is of paramount concern in that these are often difficult to reverse once established, even with large amounts of antivenom. Prolonged intubation and ventilatory support (perhaps up to a week or longer) may be required. Early diagnosis of neurotoxic symptoms and prompt and adequate dosages of antivenom are critical to avoid these complications.
The venom also contains a potent hemotoxin (procoagulants), a prothrombin activator that leads to the consumption of major coagulation factors, including fibrinogen, leading to interference with blood clotting. This causes defibrination, with nonclottable blood, putting victims at risk of major bleeding from the bite site and can lead to more serious, sometimes fatal, internal haemorrhaging, especially in the brain. Recovering from this takes many hours after venom neutralisation has been achieved with antivenom. Taipan snake procoagulants are among the most powerful snake venom procoagulants known, though mild coagulopathy has also been reported for inland taipan envenomation (Sutherland and Tibballs, 2001).
No nephrotoxins (kidney toxins) have so far been isolated from inland taipan snake venoms, but renal (kidney) impairment or acute kidney failure can occur secondary to severe rhabdomyolysis.
Taipan snake venom does contain myotoxins that cause myolysis (rhabdomyolysis, muscle damage); the urine of bite victims often turns reddish-brown as their muscles release myoglobin, which is passed through the kidneys (myoglobinuria). The kidneys are often badly damaged by filtering so much tissue debris out of the blood, and kidney failure is a common complication in serious cases of significant envenoming.
Causes of death:
Paralysis – primary, e.g., respiratory failure; secondary, e.g., pneumonia
Coagulopathy – primary, e.g., cerebral haemorrhage; secondary, e.g., kidney failure
Kidney failure – includes secondary complications such as infections
Anaphylaxis – acute allergic reaction to venom in a patient previously exposed to taipan snake venom (e.g., reptile keeper)
Cardiac complications – likely to be secondary
Antivenom
Until 1955, the only antivenom available for general distribution for Australian snakes was the monovalent (specific) tiger snake (Notechis) antivenom, which gave varying degrees of cross-protection against the bites of most other dangerous Australian snakes. Thereafter followed specific antivenom for other common snakes, among them the coastal taipan, and finally, a polyvalent (broad-spectrum) antivenom for the bites of any unidentified snake from Australia.
The coastal taipan antivenom, known as "taipan antivenom", is effective against the inland taipan venom, as well, but it is not as effective in bite victims of the inland taipan as in those of the coastal taipan.
Taipan antivenom is produced and manufactured by the Australian Reptile Park and the Commonwealth Serum Laboratories in Melbourne.
| Biology and health sciences | Snakes | Animals |
1688526 | https://en.wikipedia.org/wiki/H%20I%20region | H I region | An HI region or H I region (read H one) is a cloud in the interstellar medium composed of neutral atomic hydrogen (HI), in addition to the local abundance of helium and other elements. (H is the chemical symbol for hydrogen, and "I" is the Roman numeral. It is customary in astronomy to use the Roman numeral I for neutral atoms, II for singly-ionized—HII is H+ in other sciences—III for doubly-ionized, e.g. OIII is O++, etc.) These regions do not emit detectable visible light (except in spectral lines from elements other than hydrogen) but are observed by the 21-cm (1,420 MHz) region spectral line. This line has a very low transition probability, so it requires large amounts of hydrogen gas for it to be seen. At ionization fronts, where HI regions collide with expanding ionized gas (such as an H II region), the latter glows brighter than it otherwise would. The degree of ionization in an HI region is very small at around 10−4 (i.e. one particle in 10,000). At typical interstellar pressures in galaxies like the Milky Way, HI regions are most stable at temperatures of either below 100 K or above several thousand K; gas between these temperatures heats or cools very quickly to reach one of the stable temperature regimes. Within one of these phases, the gas is usually considered isothermal, except near an expanding H II region. Near an expanding H II region is a dense HI region, separated from the undisturbed HI region by a shock front and from the H II region by an ionization front.
Mapping
Mapping HI emissions with a radio telescope is a technique used for determining the structure of spiral galaxies. It is also used to map gravitational disruptions between galaxies. When two galaxies collide, the material is pulled out in strands, allowing astronomers to determine which way the galaxies are moving.
HI regions effectively absorb photons that are energetic enough to ionize hydrogen, which requires an energy of 13.6 electron volts. They are ubiquitous in the Milky Way galaxy, and the Lockman Hole is one of the few "windows" for clear observations of distant objects at extreme ultraviolet and soft x-ray wavelengths.
| Physical sciences | Basics_2 | Astronomy |
14997569 | https://en.wikipedia.org/wiki/Location%20of%20Earth | Location of Earth | Knowledge of the location of Earth has been shaped by 400 years of telescopic observations, and has expanded radically since the start of the 20th century. Initially, Earth was believed to be the center of the Universe,
which consisted only of those planets visible with the naked eye and an outlying sphere of fixed stars. After the acceptance of the heliocentric model in the 17th century, observations by William Herschel and others showed that the Sun lay within a vast, disc-shaped galaxy of stars. By the 20th century, observations of spiral nebulae revealed that the Milky Way galaxy was one of billions in an expanding universe, grouped into clusters and superclusters. By the end of the 20th century, the overall structure of the visible universe was becoming clearer, with superclusters forming into a vast web of filaments and voids. Superclusters, filaments and voids are the largest coherent structures in the Universe that we can observe. At still larger scales (over 1000 megaparsecs) the Universe becomes homogeneous, meaning that all its parts have on average the same density, composition and structure.
Since there is believed to be no "center" or "edge" of the Universe, there is no particular reference point with which to plot the overall location of the Earth in the universe. Because the observable universe is defined as that region of the Universe visible to terrestrial observers, Earth is, because of the constancy of the speed of light, the center of Earth's observable universe. Reference can be made to the Earth's position with respect to specific structures, which exist at various scales. It is still undetermined whether the Universe is infinite. There have been numerous hypotheses that the known universe may be only one such example within a higher multiverse; however, no direct evidence of any sort of multiverse has been observed, and some have argued that the hypothesis is not falsifiable.
Details
Earth is the third planet from the Sun with an approximate distance of , and is traveling nearly through outer space.
Table
Gallery
| Physical sciences | Solar System | Astronomy |
11353408 | https://en.wikipedia.org/wiki/Earthquake%20forecasting | Earthquake forecasting | Earthquake forecasting is a branch of the science of seismology concerned with the probabilistic assessment of general earthquake seismic hazard, including the frequency and magnitude of damaging earthquakes in a given area over years or decades. While forecasting is usually considered to be a type of prediction, earthquake forecasting is often differentiated from earthquake prediction, Earthquake forecasting estimates the likelihood of earthquakes in a specific timeframe and region, while earthquake prediction attempts to pinpoint the exact time, location, and magnitude of an impending quake, which is currently not reliably achievable.. says: "This definition has several defects which contribute to confusion and difficulty in prediction research." In addition to specification of time, location, and magnitude, Allen suggested three other requirements: 4) indication of the author's confidence in the prediction, 5) the chance of an earthquake occurring anyway as a random event, and 6) publication in a form that gives failures the same visibility as successes. define prediction (in part) "to be a formal rule where by the available space-time-seismic moment manifold of earthquake occurrence is significantly contracted ...."</ref> Both forecasting and prediction of earthquakes are distinguished from earthquake warning systems, which, upon detection of an earthquake, provide a real-time warning to regions that might be affected.
In the 1970s, scientists were optimistic that a practical method for predicting earthquakes would soon be found, but by the 1990s continuing failure led many to question whether it was even possible. Demonstrably successful predictions of large earthquakes have not occurred, and the few claims of success are controversial. Consequently, many scientific and government resources have been used for probabilistic seismic hazard estimates rather than prediction of individual earthquakes. Such estimates are used to establish building codes, insurance rate structures, awareness and preparedness programs, and public policy related to seismic events. In addition to regional earthquake forecasts, such seismic hazard calculations can take factors such as local geological conditions into account. Anticipated ground motion can then be used to guide building design criteria.
Methods for earthquake forecasting
Methods for earthquake forecasting generally look for trends or patterns that lead to an earthquake. As these trends may be complex and involve many variables, advanced statistical techniques are often needed to understand them, therefore these are sometimes called statistical methods. These approaches tend to have relatively long time periods, making them useful for earthquake forecasting.
Elastic rebound
Even the stiffest of rock is not perfectly rigid. Given a large force (such as between two immense tectonic plates moving past each other) the Earth's crust will bend or deform. According to the elastic rebound theory of , eventually the deformation (strain) becomes great enough that something breaks, usually at an existing fault. Slippage along the break (an earthquake) allows the rock on each side to rebound to a less deformed state. In the process, energy is released in various forms, including seismic waves. The cycle of tectonic force being accumulated in elastic deformation and released in a sudden rebound is then repeated. As the displacement from a single earthquake ranges from less than a meter to around 10 meters (for an M 8 quake), the demonstrated existence of large strike-slip displacements of hundreds of miles shows the existence of a long-running earthquake cycle.
Characteristic earthquakes
The most studied earthquake faults (such as the Nankai megathrust, the Wasatch fault, and the San Andreas Fault) appear to have distinct segments. The characteristic earthquake model postulates that earthquakes are generally constrained within these segments. As the lengths and other properties of the segments are fixed, earthquakes that rupture the entire fault should have similar characteristics. These include the maximum magnitude (which is limited by the length of the rupture), and the amount of accumulated strain needed to rupture the fault segment. Since continuous plate motions cause the strain to accumulate steadily, seismic activity on a given segment should be dominated by earthquakes of similar characteristics that recur at somewhat regular intervals. For a given fault segment, identifying these characteristic earthquakes and timing their recurrence rate (or conversely return period) should therefore inform us about the next rupture; this is the approach generally used in forecasting seismic hazard. Return periods are also used for forecasting other rare events, such as cyclones and floods, and assume that future frequency will be similar to observed frequency to date.
Extrapolation from the Parkfield earthquakes of 1857, 1881, 1901, 1922, 1934, and 1966 led to a forecast of an earthquake around 1988, or before 1993 at the latest (at the 95% confidence interval), based on the characteristic earthquake model. Instrumentation was put in place in hopes of detecting precursors of the anticipated earthquake. However, the forecasted earthquake did not occur until 2004. The failure of the Parkfield prediction experiment has raised doubt as to the validity of the characteristic earthquake model itself.
Seismic gaps
At the contact where two tectonic plates slip past each other, every section must eventually slip, as (in the long-term) none get left behind. But they do not all slip at the same time; different sections will be at different stages in the cycle of strain (deformation) accumulation and sudden rebound. In the seismic gap model, the "next big quake" should be expected not in the segments where recent seismicity has relieved the strain, but in the intervening gaps where the unrelieved strain is the greatest. This model has an intuitive appeal; it is used in long-term forecasting, and was the basis of a series of circum-Pacific (Pacific Rim) forecasts in 1979 and 1989–1991.
However, some underlying assumptions about seismic gaps are now known to be incorrect. A close examination suggests that "there may be no information in seismic gaps about the time of occurrence or the magnitude of the next large event in the region"; statistical tests of the circum-Pacific forecasts shows that the seismic gap model "did not forecast large earthquakes well". Another study concluded that a long quiet period did not increase earthquake potential.
Notable forecasts
UCERF3
The 2015 Uniform California Earthquake Rupture Forecast, Version 3, or UCERF3, is the latest official earthquake rupture forecast (ERF) for the state of California, superseding UCERF2. It provides authoritative estimates of the likelihood and severity of potentially damaging earthquake ruptures in the long- and near-term. Combining this with ground motion models produces estimates of the severity of ground shaking that can be expected during a given period (seismic hazard), and of the threat to the built environment (seismic risk). This information is used to inform engineering design and building codes, planning for disaster, and evaluating whether earthquake insurance premiums are sufficient for the prospective losses. A variety of hazard metrics can be calculated with UCERF3; a typical metric is the likelihood of a magnitude M 6.7 earthquake (the size of the 1994 Northridge earthquake) in the 30 years (typical life of a mortgage) since 2014.
UCERF3 was prepared by the Working Group on California Earthquake Probabilities (WGCEP), a collaboration between the United States Geological Survey (USGS), the California Geological Survey (CGS), and the Southern California Earthquake Center (SCEC), with significant funding from the California Earthquake Authority (CEA).
| Physical sciences | Seismology | Earth science |
1110017 | https://en.wikipedia.org/wiki/Dump%20truck | Dump truck | A dump truck, known also as a dumping truck, dump trailer, dumper trailer, dump lorry or dumper lorry or a dumper for short, is used for transporting materials (such as dirt, gravel, or demolition waste) for construction as well as coal. A typical dump truck is equipped with an open-box bed, which is hinged at the rear and equipped with hydraulic rams to lift the front, allowing the material in the bed to be deposited ("dumped") on the ground behind the truck at the site of delivery. In the UK, Australia, South Africa and India the term applies to off-road construction plants only and the road vehicle is known as a tip lorry, tipper lorry (UK, India), tipper truck, tip truck, tip trailer or tipper trailer or simply a tipper (Australia, New Zealand, South Africa).
History
The dump truck is thought to have been first conceived in the farms of late 19th century western Europe. Thornycroft developed a steam dust-cart in 1896 with a tipper mechanism. The first motorized dump trucks in the United States were developed by small equipment companies such as The Fruehauf Trailer Corporation, Galion Buggy Co. and Lauth-Juergens among many others around 1910. Hydraulic dump beds were introduced by Wood Hoist Co. shortly after. Such companies flourished during World War I due to massive wartime demand. August Fruehauf had obtained military contracts for his semi-trailer, invented in 1914 and later created the partner vehicle, the semi-truck for use in World War I. After the war, Fruehauf introduced hydraulics in his trailers. They offered hydraulic lift gates, hydraulic winches and a dump trailer for sales in the early 1920s. Fruehauf became the premier supplier of dump trailers and their famed "bathtub dump" was considered to be the best by heavy haulers, road and mining construction firms.
Companies like Galion Buggy Co. continued to grow after the war by manufacturing a number of express bodies and some smaller dump bodies that could be easily installed on either stock or converted (heavy-duty suspension and drivetrain) Model T chassis prior to 1920. Galion and Wood Mfg. Co. built all of the dump bodies offered by Ford on their heavy-duty AA and BB chassis during the 1930s. Galion (now Galion Godwin Truck Body Co.) is the oldest known truck body manufacturer still in operation today.
The first known Canadian dump truck was developed in Saint John, New Brunswick, when Robert T. Mawhinney attached a dump box to a flatbed truck in 1920. The lifting device was a winch attached to a cable that fed over sheave (pulley) mounted on a mast behind the cab. The cable was connected to the lower front end of the wooden dump box which was attached by a pivot at the back of the truck frame. The operator turned a crank to raise and lower the box.
From the 1930s Euclid, International-Harvester and Mack contributed to ongoing development. Mack modified its existing trucks with varying success. In 1934 Euclid became the first manufacturer in the world to successfully produce a dedicated off-highway truck.
Types
Today, virtually all dump trucks operate by hydraulics and they come in a variety of configurations each designed to accomplish a specific task in the construction material supply chain.
Standard dump truck
A standard dump truck is a truck chassis with a dump body mounted to the frame. The bed is raised by a vertical hydraulic ram mounted under the front of the body (known as a front post hoist configuration), or a horizontal hydraulic ram and lever arrangement between the frame rails (known as an underbody hoist configuration), and the back of the bed is hinged at the back of the truck. The tailgate (sometimes referred to as an end gate) can be configured to swing up on top hinges (and sometimes also to fold down on lower hinges) or it can be configured in the "High Lift Tailgate" format wherein pneumatic or hydraulic rams lift the gate open and up above the dump body. Some bodies, typically for hauling grain, have swing-out doors for entering the box and a metering gate/chute in the center for a more controlled dumping.
In the United States most standard dump trucks have one front steering axle and one (4x2 4-wheeler) or two (6x4 6-wheeler) rear axles which typically have dual wheels on each side. Tandem rear axles are almost always powered, front steering axles are also sometimes powered (4x4, 6x6). Unpowered axles are sometimes used to support extra weight. Most unpowered rear axles can be raised off the ground to minimize wear when the truck is empty or lightly loaded, and are commonly called "lift axles".
European Union heavy trucks often have two steering axles. Dump truck configurations are two, three, and four axles. The four-axle eight wheeler has two steering axles at the front and two powered axles at the rear and is limited to gross weight in most EU countries. The largest of the standard European dump trucks is commonly called a "centipede" and has seven axles. The front axle is the steering axle, the rear two axles are powered, and the remaining four are lift axles.
The shorter wheelbase of a standard dump truck often makes it more maneuverable than the higher capacity semi-trailer dump trucks.
Semi trailer end dump truck
A semi end dump is a tractor-trailer combination wherein the trailer itself contains the hydraulic hoist. In the US a typical semi end dump has a 3-axle tractor pulling a 2-axle trailer with dual tires, in the EU trailers often have 3 axles and single tires. The key advantage of a semi end dump is a large payload. A key disadvantage is that they are very unstable when raised in the dumping position limiting their use in many applications where the dumping location is uneven or off level. Some end dumps make use of an articulated arm (known as a stabilizer) below the box, between the chassis rails, to stabilize the load in the raised position.
Frame and Frameless end dump truck
Depending on the structure, semi trailer end dump truck can also be divided into frame trailer and frameless trailer.
The main difference between them is the different structure. The frame dump trailer has a large beam that runs along the bottom of the trailer to support it. The frameless dump trailer has no frame under the trailer but has ribs that go around the body for support and the top rail of the trailer serves as a suspension bridge for support.
The difference in structure also brings with it a difference in weight. Frame dump trailers are heavier. For the same length, a frame dump trailer weighs around 5 ton more than a frameless dump trailer.
Transfer dump truck
A transfer dump truck is a standard dump truck pulling a separate trailer with a movable cargo container, which can also be loaded with construction aggregate, gravel, sand, asphalt, klinkers, snow, wood chips, triple mix, etc.
The second aggregate container on the trailer ("B" box), is powered by an electric motor, a pneumatic motor or a hydraulic line. It rolls on small wheels, riding on rails from the trailer's frame into the empty main dump container ("A" box). This maximizes payload capacity without sacrificing the maneuverability of the standard dump truck. Transfer dump trucks are typically seen in the western United States due to the peculiar weight restrictions on highways there.
Another configuration is called a triple transfer train, consisting of a "B" and "C" box. These are common on Nevada and Utah Highways, but not in California. Depending on the axle arrangement, a triple transfer can haul up to with a special permit in certain American states. , a triple transfer costs a contractor about $105 an hour, while a A/B configuration costs about $85 per hour.
Transfer dump trucks typically haul between of aggregate per load, each truck is capable of 3–5 loads per day, generally speaking.
Truck and pup
A truck and pup is very similar to a transfer dump. It consists of a standard dump truck pulling a dump trailer. The pup trailer, unlike the transfer, has its own hydraulic ram and is capable of self-unloading.
Superdump truck
A super dump is a straight dump truck equipped with a trailing axle, a liftable, load-bearing axle rated as high as . Trailing behind the rear tandem, the trailing axle stretches the outer "bridge" measurement—the distance between the first and last axles—to the maximum overall length allowed. This increases the gross weight allowed under the federal bridge formula, which sets standards for truck size and weight. Depending on the vehicle length and axle configuration, Superdumps can be rated as high as GVW and carry of payload or more. When the truck is empty or ready to offload, the trailing axle toggles up off the road surface on two hydraulic arms to clear the rear of the vehicle. Truck owners call their trailing axle-equipped trucks Superdumps because they far exceed the payload, productivity, and return on investment of a conventional dump truck. The Superdump and trailing axle concept were developed by Strong Industries of Houston, Texas.
Semi trailer bottom dump truck
A semi bottom dump, bottom hopper, or belly dump is a (commonly) 3-axle tractor pulling a 2-axle trailer with a clam shell type dump gate in the belly of the trailer. The key advantage of a semi bottom dump is its ability to lay material in a windrow, a linear heap. In addition, a semi bottom dump is maneuverable in reverse, unlike the double and triple trailer configurations described below. These trailers may be found either of the windrow type shown in the photo or may be of the cross spread type, with the gate opening front to rear instead of left and right. The cross spread type gate will actually spread the cereal grains fairly and evenly from the width of the trailer. By comparison, the windrow-type gate leaves a pile in the middle. The cross spread type gate, on the other hand, tends to jam and may not work very well with coarse materials.
Double and triple trailer bottom dump truck
Double and triple bottom dumps consist of a 2-axle tractor pulling one single-axle semi-trailer and an additional full trailer (or two full trailers in the case of triples). These dump trucks allow the driver to lay material in windrows without leaving the cab or stopping the truck. The main disadvantage is the difficulty in backing double and triple units.
The specific type of dump truck used in any specific country is likely to be closely keyed to the weight and axle limitations of that jurisdiction. Rock, dirt, and other types of materials commonly hauled in trucks of this type are quite heavy, and almost any style of truck can be easily overloaded. Because of that, this type of truck is frequently configured to take advantage of local weight limitations to maximize the cargo. For example, within the United States, the maximum weight limit is throughout the country, except for specific bridges with lower limits. Individual states, in some instances, are allowed to authorize trucks up to . Most states that do so require that the trucks be very long, to spread the weight over more distance. It is in this context that double and triple bottoms are found within the United States.
Bumper Pull Dump Trailer
Bumper Pull personal and commercial Dump Trailers come in a variety of sizes from smaller 6x10 7,000 GVWR models to larger 7x16 High Side 14,000 GVWR models.
Dump trailers come with a range of options and features such as tarp kits, high side options, dump/spread/swing gates, remote control, scissor, telescop, dual or single cylinder lifts, and metal locking toolboxes. They offer the perfect solution for a variety of applications, including roofing, rock and mulch delivery, general contractors, skid steer grading, trash out, and recycling.
Side dump truck
A side dump truck (SDT) consists of a 3-axle tractor pulling a 2-axle semi-trailer. It has hydraulic rams that tilt the dump body onto its side, spilling the material to either the left or right side of the trailer. The key advantages of the side dump are that it allows rapid unloading and can carry more weight in the western United States. In addition, it is almost immune to upset (tipping over) while dumping, unlike the semi end dumps which are very prone to tipping over. It is, however, highly likely that a side dump trailer will tip over if dumping is stopped prematurely. Also, when dumping loose materials or cobble sized stone, the side dump can become stuck if the pile becomes wide enough to cover too much of the trailer's wheels. Trailers that dump at the appropriate angle (50° for example) avoid the problem of the dumped load fouling the path of the trailer wheels by dumping their loads further to the side of the truck, in some cases leaving sufficient clearance to walk between the dumped load and the trailer.
Winter service vehicles
Many winter service vehicles are based on dump trucks, to allow the placement of ballast to weigh the truck down or to hold sodium or calcium chloride salts for spreading on snow and ice-covered surfaces. Plowing is severe service and needs heavy-duty trucks.
Roll-off trucks
A Roll-off has a hoist and subframe, but no body, it carries removable containers. The container is loaded on the ground, then pulled onto the back of the truck with a winch and cable. The truck goes to the dumpsite, after it has been dumped the empty container is taken and placed to be loaded or stored. The hoist is raised and the container slides down the subframe so the rear is on the ground. The container has rollers on the rear and can be moved forward or back until the front of it is lowered onto the ground. The containers are usually open-topped boxes used for rubble and building debris, but rubbish compactor containers are also carried. A newer hook-lift system ("roller container" in the UK) does the same job, but lifts, lowers, and dumps the container with a boom arrangement instead of a cable and hoist.
Off-highway dump trucks
Off-highway dump trucks are heavy construction equipment and share little resemblance to highway dump trucks. Bigger off-highway dump trucks are used strictly off-road for mining and heavy dirt hauling jobs. There are two primary forms: rigid frame and articulating frame.
The term "dump" truck is not generally used by the mining industry, or by the manufacturers that build these machines. The more appropriate U.S. term for this strictly off-road vehicle is "haul truck" and the equivalent European term is "dumper".
Haul truck
Haul trucks are used in large surface mines and quarries. They have a rigid frame and conventional steering with drive at the rear wheel. As of late 2013, the largest ever production haul truck is the 450 metric ton BelAZ 75710, followed by the Liebherr T 282B, the Bucyrus MT6300AC and the Caterpillar 797F, which each have payload capacities of up to . The previous record holder being the Canadian-built Terex 33-19 "Titan", having held the record for over 25 years. Most large-size haul trucks employ Diesel-electric powertrains, using the Diesel engine to drive an AC alternator or DC generator that sends electric power to electric motors at each rear wheel. The Caterpillar 797 is unique for its size, as it employs a Diesel engine to power a mechanical powertrain, typical of most road-going vehicles and intermediary size haul trucks.
Other major manufacturers of haul trucks include SANY, XCMG, Hitachi, Komatsu, DAC, Terex, and BelAZ.
Articulated hauler
An articulated dumper is an all-wheel-drive, off-road dump truck. It has a hinge between the cab and the dump box but is distinct from a semi-trailer truck in that the power unit is a permanent fixture, not a separable vehicle. Steering is accomplished via hydraulic cylinders that pivot the entire tractor in relation to the trailer, rather than rack and pinion steering on the front axle as in a conventional dump truck. By this way of steering, the trailer's wheels follow the same path as the front wheels. Together with all-wheel drive and low center of gravity, it is highly adaptable to rough terrain. Major manufacturers include Volvo CE, Terex, John Deere, and Caterpillar.
U-shaped dump truck
U-shaped dump trucks, also known as tub-body trucks, is used to transport construction waste, it is made of high-strength super wear-resistant special steel plate directly bent, and has the characteristics of impact resistance, alternating stress resistance, corrosion resistance and so on.
1. Cleaner unloading
U-shaped dump truck, there is no dead angle at the corners of the cargo box, it is not easy to stick to the box when unloading, and the unloading is cleaner.
2. Lightweight
The U-shaped cargo box reduces its own weight through structural optimization. Now the most common U-shaped dump is to use high-strength plates. Under the premise of ensuring the strength of the car body, the thickness of the plate is reduced by about 20%, and the self-weight of the car is reduced by about 1 ton, which effectively improves the utilization factor of the load mass.
3. Strong carrying capacity. Using high-strength steel plate, high yield strength, better impact resistance and fatigue resistance. For users of ore transportation, it can reduce the damage of ore to the container.
4. Low center of gravity The U-shaped structure has a lower center of gravity, which makes the ride more stable, especially when cornering, and avoids spilling cargo.
5. Save tires The U-shaped cargo box can keep the cargo in the center, and the tires on both sides are more evenly stressed, which is beneficial to improve the life of the tires.
Dangers
Collisions
Dump trucks are normally built for some amount of off-road or construction site driving; as the driver is protected by the chassis and height of the driver's seat, bumpers are either placed high or omitted for added ground clearance. The disadvantage is that in a collision with a standard car, the entire motor section or luggage compartment goes under the truck. Thus, the passengers in the car could be more severely injured than would be common in a collision with another car. Several countries have made rules that new trucks should have bumpers approximately above ground in order to protect other drivers. There are also rules about how long the load or construction of the truck can go beyond the rear bumper to prevent cars that rear-end the truck from going under it.
Tipping
Another safety consideration is the leveling of the truck before unloading. If the truck is not parked on relatively horizontal ground, the sudden change of weight and balance due to lifting of the body and dumping of the material can cause the truck to slide, or even to tip over. The live bottom trailer is an approach to eliminate this danger.
Back-up accidents
Because of their size and the difficulty of maintaining visual contact with on-foot workers, dump trucks can be a threat, especially when backing up. Mirrors and back-up alarms provide some level of protection, and having a spotter working with the driver also decreases back-up injuries and fatalities.
Manufacturers
Ashok Leyland
Asia MotorWorks
Astra Veicoli Industriali
BelAZ
BEML
Case CE
Caterpillar Inc.
DAC
Daewoo
Dart (commercial vehicle)
Eicher Motors
Euclid Trucks
FAP
HEPCO
Hitachi Construction Machinery
Hitachi Construction Machinery (Europe)
Iveco
John Deere
Kamaz
Kenworth
Kioleides
Komatsu
KrAZ
Leader Trucks
Liebherr Group
Mack Trucks
Mahindra Trucks & Buses Ltd.
MAN SE
Mercedes-Benz
Navistar International
New Holland
Peterbilt
SANY
Scania AB
ST Kinetics
Tata
Tatra (company)
Terex Corporation
Volvo Construction Equipment
Volvo Trucks
XCMG
| Technology | Specific-purpose transportation | null |
1110369 | https://en.wikipedia.org/wiki/Battle%20axe | Battle axe | A battle axe (also battle-axe, battle ax, or battle-ax) is an axe specifically designed for combat. Battle axes were designed differently to utility axes, with blades more akin to cleavers than to wood axes. Many were suitable for use in one hand, while others were larger and were deployed two-handed.
Axes designed for warfare ranged in weight from just over , and in length from just over to upwards of , as in the case of the Danish axe or the sparth axe. Cleaving weapons longer than would arguably fall into the category of polearms.
Overview
Through the course of human history, commonplace objects have been pressed into service as weapons. Axes, by virtue of their ubiquity, are no exception. Besides axes designed for combat, there were many battle axes that doubled as tools. Axes could be modified into deadly projectiles as well (see the francisca for an example). Axes were often cheaper than swords and considerably more available.
Battle axes generally weigh far less than modern splitting axes, especially mauls, because they were designed to cut legs and arms rather than wood; consequently, slightly narrow slicing blades are the norm. This facilitates deep, devastating wounds. Moreover, a lighter weapon is much quicker to bring to bear in combat and manipulate for repeated strikes against an adversary.
The crescent-shaped heads of European battle axes of the Roman and post-Roman periods were usually made of wrought iron with a carbon steel edge or, as time elapsed across the many centuries of the medieval era, steel. The hardwood handles of military axes came to be reinforced with metal bands called langets, so that an enemy warrior could not cut the shaft. Some later specimens had all-metal handles.
Battle axes are particularly associated in Western popular imagination with the Vikings. Certainly, Scandinavian foot soldiers and maritime marauders employed them as a stock weapon during their heyday, which extended from the beginning of the 8th century to the end of the 11th century. They produced several varieties, including specialized throwing axes (see francisca) and "bearded" axes or "skeggox" (so named for their trailing lower blade edge which increased cleaving power and could be used to catch the edge of an opponent's shield and pull it down, leaving the shield-bearer vulnerable to a follow-up blow). Viking axes may have been wielded with one hand or two, depending on the length of the plain wooden haft. See Viking Age arms and armor.
History
Europe
Prehistory and the Ancient Mediterranean
Stone hand axes were in use in the Paleolithic period for hundreds of thousands of years. The first hafted stone axes appear to have been produced about 6000 BCE during the Mesolithic period. Technological development continued in the Neolithic period with the much wider usage of hard stones in addition to flint and chert and the widespread use of polishing to improve axe properties. The axes proved critical in wood working and became cult objects (for example, the entry for the Battle-axe people of Scandinavia, treated their axes as high-status cultural objects). Such stone axes were made from a wide variety of tough rocks such as picrite and other igneous or metamorphic rocks, and were widespread in the Neolithic period. Many axe heads found were probably used primarily as mauls to split wood beams, and as sledgehammers for construction purposes (such hammering stakes into the ground, for example).
Narrow axe heads made of cast metals were subsequently manufactured by artisans in the Middle East and then Europe during the Copper Age and the Bronze Age. The earliest specimens were socket-less.
More specifically, bronze battle-axe heads are attested in the archaeological record from ancient China and the New Kingdom of ancient Egypt. Some of them were suited for practical use as infantry weapons while others were clearly intended to be brandished as symbols of status and authority, judging by the quality of their decoration.
The epsilon axe was widely used during the Bronze Age by irregular infantry unable to afford better weapons. Its use was limited to Europe and the Middle East.
In the eastern Mediterranean Basin during the Iron Age, the double-bladed labrys axe was prevalent, and a hafted, single-bitted axe made of bronze or later iron was sometimes used as a weapon of war by the heavy infantry of ancient Greece, especially when confronted with thickly-armored opponents. The sagaris—described as either single bitted or double bitted—became associated by the Greeks with the mythological Amazons, though these were generally ceremonial axes rather than practical implements.
The Barbarian tribes that the Romans encountered north of the Alps did include iron war axes in their armories, alongside swords and spears. The Cantabri from the Iberian peninsula also used battle axes.
The Middle Ages
Battle axes were very common in Europe in the Migration Period and the subsequent Viking Age, and they famously figure on the 11th-century Bayeux Tapestry, which depicts Norman mounted knights pitted against Anglo-Saxon infantrymen. They continued to be employed throughout the rest of the Middle Ages, with significant combatants being noted axe wielders in the 12th, 13th, and 14th centuries.
King Stephen of England famously used a 'Dane axe' at the Battle of Lincoln 1141. One account says that he used it after his sword broke. Another says he used his sword only after his axe broke.
Richard the Lionheart was often recorded in Victorian times wielding a large war axe, though references are sometimes wildly exaggerated as befitted a national hero: "Long and long after he was quiet in his grave, his terrible battle-axe, with twenty English pounds of English steel in its mighty head..." – A Child's History of England by Charles Dickens. Richard is, however, recorded as using a Danish Axe at the relief of Jaffa. Geoffrey of Lusignan is another famous crusader associated with the axe.
Robert the Bruce, King of Scotland, used an axe to defeat Henry de Bohun in single combat at the start of the Battle of Bannockburn in 1314. Given that Bruce was wielding the axe on horseback, it is likely that it was a one handed horseman's axe. They enjoyed a sustained revival among heavily armored equestrian combatants in the 15th century.
In the 14th century, the use of axes is increasingly noted by Froissart in his Chronicle, which records the engagements between the kingdoms of France and England and the rise of professional and mercenary armies in the 14th century. King John II is recorded as using one at the Battle of Poitiers in 1356 and Sir James Douglas at the Battle of Otterburn in 1388. Bretons were apparently noted axe users, with noted mercenaries Bertrand du Guesclin and Olivier de Clisson both wielding axes in battle. In these instances the type of battle axe - whether a Danish axe, or the proto-pollaxe - is not recorded.
Most medieval European battle axes had a socketed head (meaning that the thicker, butt-end of the blade contained an opening into which a wooden haft was inserted), and some included langets—long strips of metal affixed to the faces of the haft to prevent it from being damaged during combat. Occasionally the cheeks of the axehead bore engraved, etched, punched, or inlaid decorative patterns. Late-period battle axes tended to be of all-metal construction.
Such medieval polearms as the halberd and the pollaxe were variants of the basic battle-axe form.
Steel plate-armor covering almost all of a knight's body, and incorporating features specifically designed to defeat axe and sword blades, become more common in the late 14th and early 15th century. Its development led to a generation of hafted weapons with points that concentrated impact, either to penetrate steel plate or to damage the joints of articulated plate. Increasingly daggers called misericords were carried which enabled a sharp point to be thrust though gaps in armour if an opponent was disabled or being grappled with. Swords styles became more diverse – from the two-handed zweihänders to more narrow thrusting instruments with sharply pointed tips, capable of penetrating any "chinks in the armour" of a fully encased opponent: for example, the estoc.
A sharp, sometimes curved pick was often fitted to the rear of the battle axe's blade to provide the user with a secondary weapon of penetration. A stabbing spike could be added, too, as a finial. Similarly, the war hammer evolved in late-medieval times with fluted or spiked heads, which would help a strike to "bite" into the armour and deliver its energy through to the wearer, rather than glance off the armor's surface. Strikes from these armour penetrating picks were not always fatal. There are many accounts of plate armored knights being struck with said weapons and while the armour was damaged, the individual underneath survived and in some cases completely unharmed.
It eventually became common for these various kinds of impact weapons to be made entirely from metal, thus doing away with reinforced wooden shafts.
A useful visual guide to high-medieval battle axes, contemporary with their employment, are the scenes of warfare depicted in the Maciejowski Bible (Morgan Bible) of c. 1250.
Battle axes also came to figure as heraldic devices on the coats of arms of several English and mainland European families.
Post-medieval axes
Battle axes were eventually phased out at the end of the 16th century as military tactics began to revolve increasingly around the use of gunpowder. However, as late as the 1640s, Prince Rupert—a Royalist general and cavalry commander during the English Civil War—is pictured carrying a battle axe, and this was not merely a decorative symbol of authority: the "short pole-axe" was adopted by Royalist cavalry officers to penetrate Roundhead troopers' helmets and cuirasses in close-quarters fighting, and it was also used by their opponents: Sir Bevil Grenville was slain by a Parliamentarian pole-axe at the Battle of Lansdowne, and Sir Richard Bulstrode was wounded by one at the Battle of Edgehill.
In Scandinavia, however, the battle axe continued in use alongside the halberd, crossbow and pole-axe until the start of the 18th century. The nature of Norwegian terrain in particular made pike and shot tactics impracticable in many cases. A law instituted in 1604 required all farmers to own weaponry to serve in the militia. The Norwegian peasant militia battle axe, much more wieldy than the pike or halberd and yet effective against mounted enemies, was a popular choice. Many such weapons were ornately decorated, and yet their functionality shows in the way that the axe head was mounted tilting upwards slightly, with a significant forward curve in the shaft, with the intent of making them more effective against armoured opponents by concentrating force onto a narrower spot.
During Napoleonic times, and later on in the 19th century, farriers in army service carried long and heavy axes as part of their kit. Although these could be used in an emergency for fighting, their primary use was logistical: the branded hooves of deceased military horses needed to be removed in order to prove that they had indeed died (and had not been stolen). Napoleon's Pioneer Corps also carried axes that were used for clearing vegetation—a practice employed by similar units in other armies.
Middle East
The tabarzin (, lit. "saddle axe" or "saddle hatchet") is the traditional battle axe of Persia. It bears one or two crescent-shaped blades. The long form of the tabar was about seven feet long, while a shorter version was about three feet long. What made the Persian axe unique is the very thin handle, which is very light and always metallic. The tabar became one of the main weapons throughout the Middle East, and was always carried at a soldier's waist not only in Persia but Egypt, and the Arab world from the time of the Crusades. Mamluk bodyguards were known as tabardiyya after the weapon. The tabarzin is sometimes carried as a symbolic weapon by wandering dervishes (Muslim ascetic worshippers).
Asia
China
Different types of battleaxes may be found in ancient China. In Chinese mythology, Xingtian (), a deity, uses a battle axe against other gods. The qi () and yue () are heavy axes. They were common in Zhou dynasty but fell out of favor with users due to the lack of mobility. The eventually became used only for ceremonial purposes and such battleaxes made of bronze and jade have been found. The dagger axe (ge) is another form used in ancient times.
Chinese battleaxes can be divided in three subgroups: Fu (), Yue () and Ge (). The distinction between a Yue and a Fu is that a Yue is, as a general rule, broader than a Fu. In the Shang dynasty the Yue was also a symbol of power, the bigger the Yue, the greater the power. There are a few rare examples of Yue with a round blade and a hole in the middle.
The Chinese Fu appeared in the Stone Age as a tool. In the Shang dynasty (–) the Fu began to be made from bronze, and began to be used as a weapon. However, the prominence of the Fu waned on the battlefield as the Zhou dynasty came to power. In the Warring States era iron axes started to appear. Up until the Han and Jin dynasty, after having lost its importance on the battle-field, the Fu once again appeared as the cavalry was used more often.
In the Sui and Tang dynasties there is evidence of the subdivision of the Fu. During the Song dynasty axes were popularized and many types of axes began to exist. The types include Phoenix Head Axes ( ), Invincible Axe ( ), Opening Mountain Axe ( ), Emei Axe ( ) and Chisel Head Axes ( ).
A well known novel from the Ming dynasty (1368–1644) knows as the Outlaws of the Marsh (or the Water Margin - Shui Hu Zhuan ) features a character known as Li Kui, the Black Whirlwind who wields two axes and fights naked.
In the Yuan and Ming dynasties, axes retained their use in the army. In the Qing dynasty new types of axes emerge among the Eight Banners Army with straight edges. The Green Standard Army among the Eight Banners used double axes weighing each, with a length of .
In modern Chinese wushu and Chinese opera there are many depictions of the axe. Many of these axes look thick and heavy, however, the axe heads are hollow.
Indian Subcontinent
The battle axe of ancient India was known as a parashu (or farasa in some dialects). Made from iron, bamboo, wood, or wootz steel, it usually measures though some are as long as . A typical parashu could have a single edge or double edge, with a hole for fixing a shaft. The haft is often tied with a leather sheet to provide a good grip. The cutting edge is invariably broad and the length of the haft could be about three to four feet. The parashu is often depicted in religious art as one of the weapons of Hindu deities such as Shiva and Durga. The sixth avatar of Lord Vishnu, Parashurama, is named after the weapon. Parashu are still used as domestic tools in Indian households, particularly in the villages, as well as being carried by certain sects of eremitic sadhu.
Philippines
The panabas (also known as nawi among some ethnic groups) is a traditional battle axe favored by the Moro and Lumad tribes of Mindanao, Philippines. It was also used as an agricultural or chopping tool. It ranges in size from and usually long and can be held with one or two hands. Hilts were often wrapped in rattan bindings or had metal collars. Due to its clean cutting capabilities it was also sometimes used as an execution weapon. It is said that the Moro warriors wielding panabas would follow the main group of warriors up front and would immediately charge in on any American survivors of the first wave of attack during the Philippine–American War.
Among the various Cordilleran peoples of the northern Philippines, another type of traditional battle axe, the head axe, was favored for headhunting raids. It was specialized for beheading enemy combatants but was also used as an agricultural tool. They were banned, along with headhunting practices, during the American colonial period of the Philippines in the early 20th century.
Sri Lanka
The keteriya was a type of battle axe that was used in ancient Sri Lanka. A keteriya consisted of a single edge and a short handle made of wood. This would allow the user to wield it with a single hand.
Vietnam
The battle axe is one of the most common type of weapons found in Vietnamese ancient cultures, particularly the Dong Son culture.
| Technology | Melee weapons | null |
1110611 | https://en.wikipedia.org/wiki/Muscle%20contraction | Muscle contraction | Muscle contraction is the activation of tension-generating sites within muscle cells. In physiology, muscle contraction does not necessarily mean muscle shortening because muscle tension can be produced without changes in muscle length, such as when holding something heavy in the same position. The termination of muscle contraction is followed by muscle relaxation, which is a return of the muscle fibers to their low tension-generating state.
For the contractions to happen, the muscle cells must rely on the change in action of two types of filaments: thin and thick filaments.
The major constituent of thin filaments is a chain formed by helical coiling of two strands of actin, and thick filaments dominantly consist of chains of the motor-protein myosin. Together, these two filaments form myofibrils - the basic functional organelles in the skeletal muscle system.
In vertebrates, skeletal muscle contractions are neurogenic as they require synaptic input from motor neurons. A single motor neuron is able to innervate multiple muscle fibers, thereby causing the fibers to contract at the same time. Once innervated, the protein filaments within each skeletal muscle fiber slide past each other to produce a contraction, which is explained by the sliding filament theory. The contraction produced can be described as a twitch, summation, or tetanus, depending on the frequency of action potentials. In skeletal muscles, muscle tension is at its greatest when the muscle is stretched to an intermediate length as described by the length-tension relationship.
Unlike skeletal muscle, the contractions of smooth and cardiac muscles are myogenic (meaning that they are initiated by the smooth or heart muscle cells themselves instead of being stimulated by an outside event such as nerve stimulation), although they can be modulated by stimuli from the autonomic nervous system. The mechanisms of contraction in these muscle tissues are similar to those in skeletal muscle tissues.
Muscle contraction can also be described in terms of two variables: length and tension. In natural movements that underlie locomotor activity, muscle contractions are multifaceted as they are able to produce changes in length and tension in a time-varying manner. Therefore, neither length nor tension is likely to remain the same in skeletal muscles that contract during locomotion. Contractions can be described as isometric if the muscle tension changes but the muscle length remains the same. In contrast, a muscle contraction is described as isotonic if muscle tension remains the same throughout the contraction. If the muscle length shortens, the contraction is concentric; if the muscle length lengthens, the contraction is eccentric.
Types
Muscle contractions can be described based on two variables: force and length. Force itself can be differentiated as either tension or load. Muscle tension is the force exerted by the muscle on an object whereas a load is the force exerted by an object on the muscle. When muscle tension changes without any corresponding changes in muscle length, the muscle contraction is described as isometric. If the muscle length changes while muscle tension remains the same, then the muscle contraction is isotonic. In an isotonic contraction, the muscle length can either shorten to produce a concentric contraction or lengthen to produce an eccentric contraction. In natural movements that underlie locomotor activity, muscle contractions are multifaceted as they are able to produce changes in length and tension in a time-varying manner. Therefore, neither length nor tension is likely to remain constant when the muscle is active during locomotor activity.
Isometric contraction
An isometric contraction of a muscle generates tension without changing length. An example can be found when the muscles of the hand and forearm grip an object; the joints of the hand do not move, but muscles generate sufficient force to prevent the object from being dropped.
Isotonic contraction
In isotonic contraction, the tension in the muscle remains constant despite a change in muscle length. This occurs when a muscle's force of contraction matches the total load on the muscle.
Concentric contraction
In concentric contraction, muscle tension is sufficient to overcome the load, and the muscle shortens as it contracts. This occurs when the force generated by the muscle exceeds the load opposing its contraction.
During a concentric contraction, a muscle is stimulated to contract according to the sliding filament theory. This occurs throughout the length of the muscle, generating a force at the origin and insertion, causing the muscle to shorten and changing the angle of the joint. In relation to the elbow, a concentric contraction of the biceps would cause the arm to bend at the elbow as the hand moved from the leg to the shoulder (a biceps curl). A concentric contraction of the triceps would change the angle of the joint in the opposite direction, straightening the arm and moving the hand towards the leg.
Eccentric contraction
In eccentric contraction, the tension generated while isometric is insufficient to overcome the external load on the muscle and the muscle fibers lengthen as they contract. Rather than working to pull a joint in the direction of the muscle contraction, the muscle acts to decelerate the joint at the end of a movement or otherwise control the repositioning of a load. This can occur involuntarily (e.g., when attempting to move a weight too heavy for the muscle to lift) or voluntarily (e.g., when the muscle is 'smoothing out' a movement or resisting gravity such as during downhill walking). Over the short-term, strength training involving both eccentric and concentric contractions appear to increase muscular strength more than training with concentric contractions alone. However, exercise-induced muscle damage is also greater during lengthening contractions.
During an eccentric contraction of the biceps muscle, the elbow starts the movement while bent and then straightens as the hand moves away from the shoulder. During an eccentric contraction of the triceps muscle, the elbow starts the movement straight and then bends as the hand moves towards the shoulder. Desmin, titin, and other z-line proteins are involved in eccentric contractions, but their mechanism is poorly understood in comparison to cross-bridge cycling in concentric contractions.
Though the muscle is doing a negative amount of mechanical work, (work is being done on the muscle), chemical energy (of fat or glucose, or temporarily stored in ATP) is nevertheless consumed, although less than would be consumed during a concentric contraction of the same force. For example, one expends more energy going up a flight of stairs than going down the same flight.
Muscles undergoing heavy eccentric loading suffer greater damage when overloaded (such as during muscle building or strength training exercise) as compared to concentric loading. When eccentric contractions are used in weight training, they are normally called negatives. During a concentric contraction, contractile muscle myofilaments of myosin and actin slide past each other, pulling the Z-lines together. During an eccentric contraction, the myofilaments slide past each other the opposite way, though the actual movement of the myosin heads during an eccentric contraction is not known. Exercise featuring a heavy eccentric load can actually support a greater weight (muscles are approximately 40% stronger during eccentric contractions than during concentric contractions) and also results in greater muscular damage and delayed onset muscle soreness one to two days after training. Exercise that incorporates both eccentric and concentric muscular contractions (i.e., involving a strong contraction and a controlled lowering of the weight) can produce greater gains in strength than concentric contractions alone. While unaccustomed heavy eccentric contractions can easily lead to overtraining, moderate training may confer protection against injury.
Eccentric contractions in movement
Eccentric contractions normally occur as a braking force in opposition to a concentric contraction to protect joints from damage. During virtually any routine movement, eccentric contractions assist in keeping motions smooth, but can also slow rapid movements such as a punch or throw. Part of training for rapid movements such as pitching during baseball involves reducing eccentric braking allowing a greater power to be developed throughout the movement.
Eccentric contractions are being researched for their ability to speed rehabilitation of weak or injured tendons. Achilles tendinitis and patellar tendonitis (also known as jumper's knee or patellar tendonosis) have been shown to benefit from high-load eccentric contractions.
Vertebrate
In vertebrate animals, there are three types of muscle tissues: skeletal, smooth, and cardiac. Skeletal muscle constitutes the majority of muscle mass in the body and is responsible for locomotor activity. Smooth muscle forms blood vessels, the gastrointestinal tract, and other areas in the body that produce sustained contractions. Cardiac muscle makes up the heart, which pumps blood. Skeletal and cardiac muscles are called striated muscle because of their striped appearance under a microscope, which is due to the highly organized alternating pattern of A bands and I bands.
Skeletal muscle
Excluding reflexes, all skeletal muscle contractions occur as a result of signals originating in the brain. The brain sends electrochemical signals through the nervous system to the motor neuron that innervates several muscle fibers. In the case of some reflexes, the signal to contract can originate in the spinal cord through a feedback loop with the grey matter. Other actions such as locomotion, breathing, and chewing have a reflex aspect to them: the contractions can be initiated either consciously or unconsciously.
Neuromuscular junction
A neuromuscular junction is a chemical synapse formed by the contact between a motor neuron and a muscle fiber. It is the site in which a motor neuron transmits a signal to a muscle fiber to initiate muscle contraction. The sequence of events that results in the depolarization of the muscle fiber at the neuromuscular junction begins when an action potential is initiated in the cell body of a motor neuron, which is then propagated by saltatory conduction along its axon toward the neuromuscular junction. Once it reaches the terminal bouton, the action potential causes a ion influx into the terminal by way of the voltage-gated calcium channels. The influx causes synaptic vesicles containing the neurotransmitter acetylcholine to fuse with the plasma membrane, releasing acetylcholine into the synaptic cleft between the motor neuron terminal and the neuromuscular junction of the skeletal muscle fiber. Acetylcholine diffuses across the synapse and binds to and activates nicotinic acetylcholine receptors on the neuromuscular junction. Activation of the nicotinic receptor opens its intrinsic sodium/potassium channel, causing sodium to rush in and potassium to trickle out. As a result, the sarcolemma reverses polarity and its voltage quickly jumps from the resting membrane potential of -90mV to as high as +75mV as sodium enters. The membrane potential then becomes hyperpolarized when potassium exits and is then adjusted back to the resting membrane potential. This rapid fluctuation is called the end-plate potential. The voltage-gated ion channels of the sarcolemma next to the end plate open in response to the end plate potential. They are sodium and potassium specific and only allow one through. This wave of ion movements creates the action potential that spreads from the motor end plate in all directions. If action potentials stop arriving, then acetylcholine ceases to be released from the terminal bouton. The remaining acetylcholine in the synaptic cleft is either degraded by active acetylcholine esterase or reabsorbed by the synaptic knob and none is left to replace the degraded acetylcholine.
Excitation–contraction coupling
Excitation–contraction coupling (ECC) is the process by which a muscular action potential in the muscle fiber causes myofibrils to contract. In skeletal muscles, excitation–contraction coupling relies on a direct coupling between two key proteins, the sarcoplasmic reticulum (SR) calcium release channel identified as the ryanodine receptor 1 (RYR1) and the voltage-gated L-type calcium channel identified as dihydropyridine receptors, (DHPRs). DHPRs are located on the sarcolemma (which includes the surface sarcolemma and the transverse tubules), while the RyRs reside across the SR membrane. The close apposition of a transverse tubule and two SR regions containing RyRs is described as a triad and is predominantly where excitation–contraction coupling takes place.
Excitation–contraction coupling (ECC) occurs when depolarization of skeletal muscles (usually through neural innervation) results in a muscle action potential. This action potential spreads across the muscle's surface and into the muscle fiber's network of T-tubules, depolarizing the inner portion of the muscle fiber. This activates dihydropyridine receptors in the terminal cisternae, which are in close proximity to ryanodine receptors in the adjacent sarcoplasmic reticulum. The activated dihydropyridine receptors physically interact with ryanodine receptors to activate them via foot processes (involving conformational changes that allosterically activates the ryanodine receptors). As ryanodine receptors open, Ca2+ is released from the sarcoplasmic reticulum into the local junctional space and diffuses into the bulk cytoplasm to cause a calcium spark. The action potential creates a near synchronous activation of thousands of calcium sparks and causes a cell-wide increase in calcium giving rise to the upstroke of the calcium transient. The Ca2+ released into the cytosol binds to Troponin C by the actin filaments. This bond allows the actin filaments to perform cross-bridge cycling, producing force and, in some situations, motion.
When the desired motion is accomplished, relaxation can be achieved quickly through numerous pathways. Relaxation is quickly achieved through a Ca2+ buffer with various cytoplasmic proteins binding to Ca2+ with very high affinity. These cytoplasmic proteins allow for quick relaxation in fast twitch muscles. Although slower, the sarco/endoplasmic reticulum calcium-ATPase (SERCA) actively pumps Ca2+ back into the sarcoplasmic reticulum, resulting in a permanent relaxation until the next action potential arrives.
Mitochondria also participate in Ca2+ reuptake, ultimately delivering their gathered Ca2+ to SERCA for storage in the sarcoplasmic reticulum. A few of the relaxation mechanisms (NCX, Ca2+ pumps and Ca2+ leak channels) move Ca2+ completely out of the cells as well. As Ca2+ concentration declines to resting levels, Ca2+ releases from Troponin C, disallowing cross bridge-cycling, causing the force to decline and relaxation to occur. Once relaxation has fully occurred, the muscle is able to contract again, thus fully resetting the cycle.
Sliding filament theory
The sliding filament theory describes a process used by muscles to contract. It is a cycle of repetitive events that cause a thin filament to slide over a thick filament and generate tension in the muscle. It was independently developed by Andrew Huxley and Rolf Niedergerke and by Hugh Huxley and Jean Hanson in 1954. Physiologically, this contraction is not uniform across the sarcomere; the central position of the thick filaments becomes unstable and can shift during contraction but this is countered by the actions of the elastic myofilament of titin. This fine myofilament maintains uniform tension across the sarcomere by pulling the thick filament into a central position.
Cross-bridge cycle
Cross-bridge cycling is a sequence of molecular events that underlies the sliding filament theory. A cross-bridge is a myosin projection, consisting of two myosin heads, that extends from the thick filaments. Each myosin head has two binding sites: one for adenosine triphosphate (ATP) and another for actin. The binding of ATP to a myosin head detaches myosin from actin, thereby allowing myosin to bind to another actin molecule. Once attached, the ATP is hydrolyzed by myosin, which uses the released energy to move into the "cocked position" whereby it binds weakly to a part of the actin binding site. The remainder of the actin binding site is blocked by tropomyosin. With the ATP hydrolyzed, the cocked myosin head now contains adenosine diphosphate (ADP) + Pi. Two ions bind to troponin C on the actin filaments. The troponin- complex causes tropomyosin to slide over and unblock the remainder of the actin binding site. Unblocking the rest of the actin binding sites allows the two myosin heads to close and myosin to bind strongly to actin. The myosin head then releases the inorganic phosphate and initiates a power stroke, which generates a force of 2 pN. The power stroke moves the actin filament inwards, thereby shortening the sarcomere. Myosin then releases ADP but still remains tightly bound to actin. At the end of the power stroke, ADP is released from the myosin head, leaving myosin attached to actin in a rigor state until another ATP binds to myosin. A lack of ATP would result in the rigor state characteristic of rigor mortis. Once another ATP binds to myosin, the myosin head will again detach from actin and another cross-bridge cycle occurs.
Cross-bridge cycling is able to continue as long as there are sufficient amounts of ATP and in the cytoplasm. Termination of cross-bridge cycling can occur when is actively pumped back into the sarcoplasmic reticulum. When is no longer present on the thin filament, the tropomyosin changes conformation back to its previous state so as to block the binding sites again. The myosin ceases binding to the thin filament, and the muscle relaxes. The ions leave the troponin molecule to maintain the ion concentration in the sarcoplasm. The active pumping of ions into the sarcoplasmic reticulum creates a deficiency in the fluid around the myofibrils. This causes the removal of ions from the troponin. Thus, the tropomyosin-troponin complex again covers the binding sites on the actin filaments and contraction ceases.
Gradation of skeletal muscle contractions
The strength of skeletal muscle contractions can be broadly separated into twitch, summation, and tetanus. A twitch is a single contraction and relaxation cycle produced by an action potential within the muscle fiber itself. The time between a stimulus to the motor nerve and the subsequent contraction of the innervated muscle is called the latent period, which usually takes about 10 ms and is caused by the time taken for nerve action potential to propagate, the time for chemical transmission at the neuromuscular junction, then the subsequent steps in excitation-contraction coupling.
If another muscle action potential were to be produced before the complete relaxation of a muscle twitch, then the next twitch will simply sum onto the previous twitch, thereby producing a summation. Summation can be achieved in two ways: frequency summation and multiple fiber summation. In frequency summation, the force exerted by the skeletal muscle is controlled by varying the frequency at which action potentials are sent to muscle fibers. Action potentials do not arrive at muscles synchronously, and, during a contraction, some fraction of the fibers in the muscle will be firing at any given time. In a typical circumstance, when humans are exerting their muscles as hard as they are consciously able, roughly one-third of the fibers in each of those muscles will fire at once, though this ratio can be affected by various physiological and psychological factors (including Golgi tendon organs and Renshaw cells). This 'low' level of contraction is a protective mechanism to prevent avulsion of the tendon—the force generated by a 95% contraction of all fibers is sufficient to damage the body. In multiple fiber summation, if the central nervous system sends a weak signal to contract a muscle, the smaller motor units, being more excitable than the larger ones, are stimulated first. As the strength of the signal increases, more motor units are excited in addition to larger ones, with the largest motor units having as much as 50 times the contractile strength as the smaller ones. As more and larger motor units are activated, the force of muscle contraction becomes progressively stronger. A concept known as the size principle, allows for a gradation of muscle force during weak contraction to occur in small steps, which then become progressively larger when greater amounts of force are required.
Finally, if the frequency of muscle action potentials increases such that the muscle contraction reaches its peak force and plateaus at this level, then the contraction is a tetanus.
Length-tension relationship
Length-tension relationship relates the strength of an isometric contraction to the length of the muscle at which the contraction occurs. Muscles operate with greatest active tension when close to an ideal length (often their resting length). When stretched or shortened beyond this (whether due to the action of the muscle itself or by an outside force), the maximum active tension generated decreases. This decrease is minimal for small deviations, but the tension drops off rapidly as the length deviates further from the ideal. Due to the presence of elastic proteins within a muscle cell (such as titin) and extracellular matrix, as the muscle is stretched beyond a given length, there is an entirely passive tension, which opposes lengthening. Combined, there is a strong resistance to lengthening an active muscle far beyond the peak of active tension.
Force-velocity relationships
Force–velocity relationship relates the speed at which a muscle changes its length (usually regulated by external forces, such as load or other muscles) to the amount of force that it generates. Force declines in a hyperbolic fashion relative to the isometric force as the shortening velocity increases, eventually reaching zero at some maximum velocity. The reverse holds true for when the muscle is stretched – force increases above isometric maximum, until finally reaching an absolute maximum. This intrinsic property of active muscle tissue plays a role in the active damping of joints that are actuated by simultaneously active opposing muscles. In such cases, the force-velocity profile enhances the force produced by the lengthening muscle at the expense of the shortening muscle. This favoring of whichever muscle returns the joint to equilibrium effectively increases the damping of the joint. Moreover, the strength of the damping increases with muscle force. The motor system can thus actively control joint damping via the simultaneous contraction (co-contraction) of opposing muscle groups.
Smooth muscle
Smooth muscles can be divided into two subgroups: single-unit and multiunit. Single-unit smooth muscle cells can be found in the gut and blood vessels. Because these cells are linked together by gap junctions, they are able to contract as a functional syncytium. Single-unit smooth muscle cells contract myogenically, which can be modulated by the autonomic nervous system.
Unlike single-unit smooth muscle cells, multiunit smooth muscle cells are found in the muscle of the eye and in the base of hair follicles. Multiunit smooth muscle cells contract by being separately stimulated by nerves of the autonomic nervous system. As such, they allow for fine control and gradual responses, much like motor unit recruitment in skeletal muscle.
Mechanisms of smooth muscle contraction
The contractile activity of smooth muscle cells can be tonic (sustained) or phasic (transient) and is influenced by multiple inputs such as spontaneous electrical activity, neural and hormonal inputs, local changes in chemical composition, and stretch. This is in contrast to the contractile activity of skeletal muscle cells, which relies on a single neural input. Some types of smooth muscle cells are able to generate their own action potentials spontaneously, which usually occur following a pacemaker potential or a slow wave potential. These action potentials are generated by the influx of extracellular , and not . Like skeletal muscles, cytosolic ions are also required for crossbridge cycling in smooth muscle cells.
The two sources for cytosolic in smooth muscle cells are the extracellular entering through calcium channels and the ions that are released from the sarcoplasmic reticulum. The elevation of cytosolic results in more binding to calmodulin, which then binds and activates myosin light-chain kinase. The calcium-calmodulin-myosin light-chain kinase complex phosphorylates myosin on the 20 kilodalton (kDa) myosin light chains on amino acid residue-serine 19, enabling the molecular interaction of myosin and actin, and initiating contraction and activating the myosin ATPase. Unlike skeletal muscle cells, smooth muscle cells lack troponin, even though they contain the thin filament protein tropomyosin and other notable proteins – caldesmon and calponin. Thus, smooth muscle contractions are initiated by the -activated phosphorylation of myosin rather than binding to the troponin complex that regulates myosin binding sites on actin like in skeletal and cardiac muscles.
Termination of crossbridge cycling (and leaving the muscle in latch-state) occurs when myosin light chain phosphatase removes the phosphate groups from the myosin heads. Phosphorylation of the 20 kDa myosin light chains correlates well with the shortening velocity of smooth muscle. During this period, there is a rapid burst of energy use as measured by oxygen consumption. Within a few minutes of initiation, the calcium level markedly decreases, the 20 kDa myosin light chains' phosphorylation decreases, and energy use decreases; however, force in tonic smooth muscle is maintained. During contraction of muscle, rapidly cycling crossbridges form between activated actin and phosphorylated myosin, generating force. It is hypothesized that the maintenance of force results from dephosphorylated "latch-bridges" that slowly cycle and maintain force. A number of kinases such as rho kinase, DAPK3, and protein kinase C are believed to participate in the sustained phase of contraction, and flux may be significant.
Neuromodulation
Although smooth muscle contractions are myogenic, the rate and strength of their contractions can be modulated by the autonomic nervous system. Postganglionic nerve fibers of parasympathetic nervous system release the neurotransmitter acetylcholine, which binds to muscarinic acetylcholine receptors (mAChRs) on smooth muscle cells. These receptors are metabotropic, or G-protein coupled receptors that initiate a second messenger cascade. Conversely, postganglionic nerve fibers of the sympathetic nervous system release the neurotransmitters epinephrine and norepinephrine, which bind to adrenergic receptors that are also metabotropic. The exact effects on the smooth muscle depend on the specific characteristics of the receptor activated—both parasympathetic input and sympathetic input can be either excitatory (contractile) or inhibitory (relaxing).
Cardiac muscle
There are two types of cardiac muscle cells: autorhythmic and contractile. Autorhythmic cells do not contract, but instead set the pace of contraction for other cardiac muscle cells, which can be modulated by the autonomic nervous system. In contrast, contractile muscle cells (cardiomyocytes) constitute the majority of the heart muscle and are able to contract.
Excitation-contraction coupling
In both skeletal and cardiac muscle excitation-contraction (E-C) coupling, depolarization conduction and Ca2+ release processes occur. However, though the proteins involved are similar, they are distinct in structure and regulation. The dihydropyridine receptors (DHPRs) are encoded by different genes, and the ryanodine receptors (RyRs) are distinct isoforms. Besides, DHPR contacts with RyR1 (main RyR isoform in skeletal muscle) to regulate Ca2+ release in skeletal muscle, while the L-type calcium channel (DHPR on cardiac myocytes) and RyR2 (main RyR isoform in cardiac muscle) are not physically coupled in cardiac muscle, but face with each other by a junctional coupling.
Unlike skeletal muscle, E-C coupling in cardiac muscle is thought to depend primarily on a mechanism called calcium-induced calcium release, which is based on the junctional structure between T-tubule and sarcoplasmic reticulum. Junctophilin-2 (JPH2) is essential to maintain this structure, as well as the integrity of T-tubule. Another protein, receptor accessory protein 5 (REEP5), functions to keep the normal morphology of junctional SR. Defects of junctional coupling can result from deficiencies of either of the two proteins. During the process of calcium-induced calcium release, RyR2s are activated by a calcium trigger, which is brought about by the flow of Ca2+ through the L-type calcium channels. After this, cardiac muscle tends to exhibit diad structures, rather than triads.
Excitation-contraction coupling in cardiac muscle cells occurs when an action potential is initiated by pacemaker cells in the sinoatrial node or atrioventricular node and conducted to all cells in the heart via gap junctions. The action potential travels along the surface membrane into T-tubules (the latter are not seen in all cardiac cell types) and the depolarisation causes extracellular to enter the cell via L-type calcium channels and possibly sodium-calcium exchanger (NCX) during the early part of the plateau phase. Although this Ca2+ influx only count for about 10% of the Ca2+ needed for activation, it is relatively larger than that of skeletal muscle. This influx causes a small local increase in intracellular . The increase of intracellular is detected by RyR2 in the membrane of the sarcoplasmic reticulum, which releases in a positive feedback physiological response. This positive feedback is known as calcium-induced calcium release and gives rise to calcium sparks ( sparks). The spatial and temporal summation of ~30,000 sparks gives a cell-wide increase in cytoplasmic calcium concentration. The increase in cytosolic calcium following the flow of calcium through the cell membrane and sarcoplasmic reticulum is moderated by calcium buffers, which bind a large proportion of intracellular calcium. As a result, a large increase in total calcium leads to a relatively small rise in free .
The cytoplasmic calcium binds to Troponin C, moving the tropomyosin complex off the actin binding site allowing the myosin head to bind to the actin filament. From this point on, the contractile mechanism is essentially the same as for skeletal muscle (above). Briefly, using ATP hydrolysis, the myosin head pulls the actin filament toward the centre of the sarcomere.
Following systole, intracellular calcium is taken up by the sarco/endoplasmic reticulum ATPase (SERCA) pump back into the sarcoplasmic reticulum ready for the next cycle to begin. Calcium is also ejected from the cell mainly by the sodium-calcium exchanger (NCX) and, to a lesser extent, a plasma membrane calcium ATPase. Some calcium is also taken up by the mitochondria. An enzyme, phospholamban, serves as a brake for SERCA. At low heart rates, phospholamban is active and slows down the activity of the ATPase so that does not have to leave the cell entirely. At high heart rates, phospholamban is phosphorylated and deactivated thus taking most from the cytoplasm back into the sarcoplasmic reticulum. Once again, calcium buffers moderate this fall in concentration, permitting a relatively small decrease in free concentration in response to a large change in total calcium. The falling concentration allows the troponin complex to dissociate from the actin filament thereby ending contraction. The heart relaxes, allowing the ventricles to fill with blood and begin the cardiac cycle again.
Invertebrate
Circular and longitudinal muscles
In annelids such as earthworms and leeches, circular and longitudinal muscles cells form the body wall of these animals and are responsible for their movement. In an earthworm that is moving through a soil, for example, contractions of circular and longitudinal muscles occur reciprocally while the coelomic fluid serves as a hydroskeleton by maintaining turgidity of the earthworm. When the circular muscles in the anterior segments contract, the anterior portion of animal's body begins to constrict radially, which pushes the incompressible coelomic fluid forward and increasing the length of the animal. As a result, the front end of the animal moves forward. As the front end of the earthworm becomes anchored and the circular muscles in the anterior segments become relaxed, a wave of longitudinal muscle contractions passes backwards, which pulls the rest of animal's trailing body forward. These alternating waves of circular and longitudinal contractions is called peristalsis, which underlies the creeping movement of earthworms.
Obliquely striated muscles
Invertebrates such as annelids, mollusks, and nematodes, possess obliquely striated muscles, which contain bands of thick and thin filaments that are arranged helically rather than transversely, like in vertebrate skeletal or cardiac muscles. In bivalves, the obliquely striated muscles can maintain tension over long periods without using too much energy. Bivalves use these muscles to keep their shells closed.
Asynchronous muscles
Advanced insects such as wasps, flies, bees, and beetles possess asynchronous muscles that constitute the flight muscles in these animals. These flight muscles are often called fibrillar muscles because they contain myofibrils that are thick and conspicuous. A remarkable feature of these muscles is that they do not require stimulation for each muscle contraction. Hence, they are called asynchronous muscles because the number of contractions in these muscles do not correspond (or synchronize) with the number of action potentials. For example, a wing muscle of a tethered fly may receive action potentials at a frequency of 3 Hz but it is able to beat at a frequency of 120 Hz. The high frequency beating is made possible because the muscles are connected to a resonant system, which is driven to a natural frequency of vibration.
History
In 1780, Luigi Galvani discovered that the muscles of dead frogs' legs twitched when struck by an electrical spark. This was one of the first forays into the study of bioelectricity, a field that still studies the electrical patterns and signals in tissues such as nerves and muscles.
In 1952, the term excitation–contraction coupling was coined to describe the physiological process of converting an electrical stimulus to a mechanical response. This process is fundamental to muscle physiology, whereby the electrical stimulus is usually an action potential and the mechanical response is contraction. Excitation–contraction coupling can be dysregulated in many diseases. Though excitation–contraction coupling has been known for over half a century, it is still an active area of biomedical research. The general scheme is that an action potential arrives to depolarize the cell membrane. By mechanisms specific to the muscle type, this depolarization results in an increase in cytosolic calcium that is called a calcium transient. This increase in calcium activates calcium-sensitive contractile proteins that then use ATP to cause cell shortening.
The mechanism for muscle contraction evaded scientists for years and requires continued research and updating. The sliding filament theory was independently developed by Andrew F. Huxley and Rolf Niedergerke and by Hugh Huxley and Jean Hanson. Their findings were published as two consecutive papers published in the 22 May 1954 issue of Nature under the common theme "Structural Changes in Muscle During Contraction".
| Biology and health sciences | Basics_3 | null |
1111506 | https://en.wikipedia.org/wiki/Lepidodendron | Lepidodendron | Lepidodendron is an extinct genus of primitive lycopodian vascular plants belonging the order Lepidodendrales. It is well preserved and common in the fossil record. Like other Lepidodendrales, species of Lepidodendron grew as large-tree-like plants in wetland coal forest environments. They sometimes reached heights of , and the trunks were often over in diameter. They are often known as "scale trees", due to their bark having been covered in diamond shaped leaf-bases, from which leaves grew during earlier stages of growth. However, they are correctly defined as arborescent lycophytes. They thrived during the Carboniferous Period (358.9 to 298.9 million years ago), and persisted until the end of the Permian around 252 million years ago. Sometimes erroneously called "giant club mosses", the genus was actually more closely related to modern quillworts than to modern club mosses. In the form classification system used in paleobotany, Lepidodendron is both used for the whole plant as well as specifically the stems and leaves.
Etymology
The name Lepidodendron comes from the Greek λεπίς , scale, and δένδρον dendron, tree.
Description and biology
Overview
Lepidodendron species were comparable in size to modern trees. The plants had tapering trunks as wide as at their base that rose to about and even , arising from an underground system of horizontally spreading branches that were covered with many rootlets. Though the height of the lycopsids make the plants similar to modern trees, the constant dichotomy of branches created a habit that contrasts with that of modern trees. At the ends of branches were oval-shaped strobili called Lepidostrobus that had a similar shape to modern cones of a spruce or fir.
Stem
The stem of the lycopsids had a unifacial vascular cambium, contrasting with the bifacial vascular cambium of modern trees. Though the bifacial cambium of modern trees produces both secondary phloem and xylem, the unifacial cambium of Lepidodendron lycopsid produced only secondary xylem. As the lycopods aged, the wood produced by the unifacial cambium decreased towards the top of the plant such that terminal twigs resembled young Lepidodendron stems. Compared to modern trees, the stems and branches of the lycopsids contained little wood with the majority of mature stems consisting of a massive cortical meristem. The nearly-uniform growth of this cortical tissue indicates no difference in growth during changing seasons, and the absence of dormant buds further indicates the lack of seasonality in Lepidodendron species. The outermost cortex of oldest stems developed into the bark-like lycopodiopsid periderm. The bark of the lycopsid was somewhat similar to that of Picea species, as leaf scars formed peg-like projections that stretched and tore as the bark stretched. To resist the bending force of wind, Lepidodendron depended on their outer bark rather than their vascular tissues, as compared to modern trees that rely mostly on their central mass of wood.
Leaves
The leaves of the lycopsid were needle-like and were densely spiraled about young shoots, each possessing only a single vein. The leaves were similar to those of a fir in some species and similar to those of Pinus roxburghii in others, though in general the leaves of Lepidodendron species are indistinguishable from those of Sigillaria species. The decurrent leaves formed a cylindrical shell around branches. The leaves were only present on thin and young branches, indicating that, though the lycopsid were evergreen, they did not retain their needles for as long as modern conifers.
The leaf-cushions were fusiform and elongated, growing at most to a length of and a width of . The middle of leaf-cushions were smooth, where leaf scars were created when an abscission layer cut a leaf from its base. Each leaf scar was composed of a central circular or triangular scar and two lateral scars that were smaller and oval-shaped. This central scar marks where the main vascular bundle of the leaf connected to the vascular system of the stem. This xylem bundle was composed only of primary trachea. The two outer scars mark the forked branches of a strand of vascular tissue that passed from the cortex of the stem into the leaf. This forked strand is sometimes referred to as the "parichnos". Surrounding this strand were parenchyma cells and occasionally thick-walled elements. Surrounding both conducting tissues was a broad sheath of transfusion tracheids. Below the leaf scar the leaf-cushion tapered to a basal position. In this tapering area, circular impressions with fine pits were present. These impressions were continuous with the parichnos scars near the top of the tapering portion. This is because the impressions are formed by aerenchyma tissue that developed in closely with the parichnos. Above the leaf scar was a deep triangular impression known as the "ligular pit" for its similarities to the ligule of Isoetes. In some leaf-cushions a second depression was present above the ligular pit. Though its purpose is unclear, it has been suggested that the depression may mark the position of a sporangium. As the branch of a Lepidodendron lycopsid grew the leaf-cushion only grew to a certain extent, past which the leaf-cushion stretched. This stretching widened the groove that separated the leaf-cushions, creating a broad, flat channel.
Underground Structures
The underground structures of Lepidodendron and similar lycopsid species known from the fossil record including Sigillaria are assigned to the form taxon, Stigmaria. The rootlets were dichotomously branched from the rhizomes similar to Isoetes. These rhizomorphic axes were shoot-like, and dichotomous branching of the rootlets structured the stigmarian systems. Rootlet scars can be seen from Stigmaria fossils where the root hairs used to be attached. Hyphae are occasionally present in the tissues of Lepidodendron lycopsids, indicating the presence of mycorrhizal associations.
Decay
Different fossil genera have been described to name the various levels of decay in Lepidodendron bark fossils. The name Bergeria describes stems that have lost their epidermises, Aspidiariu is used when cushions have been removed by deep decay, and Knorria is used when the leaf cushions and the majority of cortical tissues has decayed, with a shallow "fluted" surface remaining. However, it has been suggested that these are more likely growth forms than preserved bark types, as entire fossilized trunks have been discovered with dissimilar forms; if decay is assumed to be constant throughout the trunk, then different forms indicate growth rather than levels of decay. It is likely that the trunk of Lepidodendron lycopsids were subject to the growth forms Knorria, Aspidiaria, and Bergeria progressing up the trunk, respectively.
Growth and reproduction
During the early stages of growth, Lepidodendron grew as single, unbranched trunk, with leaves growing out of the scale leaf bases (cushions). Towards the end of the lycopod growth, the leaves on the lower part of the trunk were shed, and in Lepidodendron, the upper part of the trunk dichotomously branched into a crown. The rate of growth of arborescent lycophytes is disputed, some authors contended that they had a rapid life cycle, growing to their maximum size and dying in only 10 to 15 years, while other authors argue that these growth rates were overestimated. Rather than reproduce with seeds, Lepidodendron lycopsids reproduced with spores. The spores were stored in sporangia situated on fertile stems that grew on or near the main trunk. The fertile stems grew together in cone-like structures that clustered at the tips of branches.
Distribution
The lack of growth rings and dormant buds indicates no seasonal growth patterns, and modern plants with similar characteristics tend to grow in tropical conditions. However, Lepidodendron species were distributed throughout subtropical regions. The lycopsid inhabited an extensive area compared to tropical flora of the same time period, with lycopods growing as far north as Spitsbergen and as far south as South America, in a latitudinal range of 120°.
Extinction
In Euramerica, Lepidodendron became extinct at the end of the Carboniferous, as part of a broader pattern of ecological change, including the increasing dominance of seed plants in lowland wetland forests, and increasingly arid-adapted vegetation across western Pangea. However, in the Cathaysia region comprising what is now China, wet tropical environmental conditions continued to prevail, with Lepidodendron (in its broad sense) only becoming extinct around the end of the Permian, around 252 million years ago, as a result of the extreme environmental disturbance caused by the Permian-Triassic extinction event.
Gallery
| Biology and health sciences | Lycophytes | Plants |
1111581 | https://en.wikipedia.org/wiki/Reaction%20%28physics%29 | Reaction (physics) | As described by the third of Newton's laws of motion of classical mechanics, all forces occur in pairs such that if one object exerts a force on another object, then the second object exerts an equal and opposite reaction force on the first. The third law is also more generally stated as: "To every action there is always opposed an equal reaction: or the mutual actions of two bodies upon each other are always equal, and directed to contrary parts." The attribution of which of the two forces is the action and which is the reaction is arbitrary. Either of the two can be considered the action, while the other is its associated reaction.
Examples
Interaction with ground
When something is exerting force on the ground, the ground will push back with equal force in the opposite direction. In certain fields of applied physics, such as biomechanics, this force by the ground is called 'ground reaction force'; the force by the object on the ground is viewed as the 'action'.
When someone wants to jump, he or she exerts additional downward force on the ground ('action'). Simultaneously, the ground exerts upward force on the person ('reaction'). If this upward force is greater than the person's weight, this will result in upward acceleration. When these forces are perpendicular to the ground, they are also called a normal force.
Likewise, the spinning wheels of a vehicle attempt to slide backward across the ground. If the ground is not too slippery, this results in a pair of friction forces: the 'action' by the wheel on the ground in backward direction, and the 'reaction' by the ground on the wheel in forward direction. This forward force propels the vehicle.
Gravitational forces
The Earth, among other planets, orbits the Sun because the Sun exerts a gravitational pull that acts as a centripetal force, holding the Earth to it, which would otherwise go shooting off into space. If the Sun's pull is considered an action, then Earth simultaneously exerts a reaction as a gravitational pull on the Sun. Earth's pull has the same amplitude as the Sun but in the opposite direction. Since the Sun's mass is so much larger than Earth's, the Sun does not generally appear to react to the pull of Earth, but in fact it does, as demonstrated in the animation (not to precise scale). A correct way of describing the combined motion of both objects (ignoring all other celestial bodies for the moment) is to say that they both orbit around the center of mass, referred to in astronomy as the barycenter, of the combined system.
Supported mass
Any mass on earth is pulled down by the gravitational force of the earth; this force is also called its weight. The corresponding 'reaction' is the gravitational force that mass exerts on the planet.
If the object is supported so that it remains at rest, for instance by a cable from which it is hanging, or by a surface underneath, or by a liquid on which it is floating, there is also a support force in upward direction (tension force, normal force, buoyant force, respectively). This support force is an 'equal and opposite' force; we know this not because of Newton's third law, but because the object remains at rest, so that the forces must be balanced.
To this support force there is also a 'reaction': the object pulls down on the supporting cable, or pushes down on the supporting surface or liquid. In this case, there are therefore four forces of equal magnitude:
F1. gravitational force by earth on object (downward)
F2. gravitational force by object on earth (upward)
F3. force by support on object (upward)
F4. force by object on support (downward)
Forces F1 and F2 are equal, due to Newton's third law; the same is true for forces F3 and F4.
Forces F1 and F3 are equal if and only if the object is in equilibrium, and no other forces are applied. (This has nothing to do with Newton's third law.)
Mass on a spring
If a mass is hanging from a spring, the same considerations apply as before. However, if this system is then perturbed (e.g., the mass is given a slight kick upwards or downwards, say), the mass starts to oscillate up and down. Because of these accelerations (and subsequent decelerations), we conclude from Newton's second law that a net force is responsible for the observed change in velocity. The gravitational force pulling down on the mass is no longer equal to the upward elastic force of the spring. In the terminology of the previous section, F1 and F3 are no longer equal.
However, it is still true that F1 = F2 and F3 = F4, as this is required by Newton's third law.
Causal misinterpretation
The terms 'action' and 'reaction' have the misleading suggestion of causality, as if the 'action' is the cause and 'reaction' is the effect. It is therefore easy to think of the second force as being there because of the first, and even happening some time after the first. This is incorrect; the forces are perfectly simultaneous, and are there for the same reason.
When the forces are caused by a person's volition (e.g. a soccer player kicks a ball), this volitional cause often leads to an asymmetric interpretation, where the force by the player on the ball is considered the 'action' and the force by the ball on the player, the 'reaction'. But physically, the situation is symmetric. The forces on ball and player are both explained by their nearness, which results in a pair of contact forces (ultimately due to electric repulsion). That this nearness is caused by a decision of the player has no bearing on the physical analysis. As far as the physics is concerned, the labels 'action' and 'reaction' can be flipped.
'Equal and opposite'
One problem frequently observed by physics educators is that students tend to apply Newton's third law to pairs of 'equal and opposite' forces acting on the same object.
This is incorrect; the third law refers to forces on two different objects. In contrast, a book lying on a table is subject to a downward gravitational force (exerted by the earth) and to an upward normal force by the table, both forces acting on the same book. Since the book is not accelerating, these forces must be exactly balanced, according to Newton's second law. They are therefore 'equal and opposite', yet they are acting on the same object, hence they are not action-reaction forces in the sense of Newton's third law. The actual action-reaction forces in the sense of Newton's third law are the weight of the book (the attraction of the Earth on the book) and the book's upward gravitational force on the earth. The book also pushes down on the table and the table pushes upwards on the book.
Moreover, the forces acting on the book are not always equally strong; they will be different if the book is pushed down by a third force, or if the table is slanted, or if the table-and-book system is in an accelerating elevator. The case of any number of forces acting on the same object is covered by considering the sum of all forces.
A possible cause of this problem is that the third law is often stated in an abbreviated form: For every action there is an equal and opposite reaction, without the details, namely that these forces act on two different objects. Moreover, there is a causal connection between the weight of something and the normal force: if an object had no weight, it would not experience support force from the table, and the weight dictates how strong the support force will be. This causal relationship is not due to the third law but to other physical relations in the system.
Centripetal and centrifugal force
Another common mistake is to state that "the centrifugal force that an object experiences is the reaction to the centripetal force on that object."
If an object were simultaneously subject to both a centripetal force and an equal and opposite centrifugal force, the resultant force would vanish and the object could not experience a circular motion. The centrifugal force is sometimes called a fictitious force or pseudo force, to underscore the fact that such a force only appears when calculations or measurements are conducted in non-inertial reference frames.
| Physical sciences | Classical mechanics | Physics |
1111758 | https://en.wikipedia.org/wiki/Antennae%20Galaxies | Antennae Galaxies | The Antennae Galaxies (also known as NGC 4038/NGC 4039 or Caldwell 60/Caldwell 61) are a pair of interacting galaxies in the constellation Corvus. They are currently going through a starburst phase, in which the collision of clouds of gas and dust, with entangled magnetic fields, causes rapid star formation. They were discovered by William Herschel in 1785.
General information
The Antennae Galaxies are undergoing a galactic collision. Located in the NGC 4038 group with five other galaxies, these two galaxies are known as the Antennae Galaxies because the two long tails of stars, gas and dust ejected from the galaxies as a result of the collision resemble an insect's antennae.
The nuclei of the two galaxies are joining to become one giant galaxy. Most galaxies probably undergo at least one significant collision in their lifetimes. This is likely the future of our Milky Way when it collides with the Andromeda Galaxy. This collision and merger sequence (the Toomre sequence) for galaxy evolution was developed in part by successfully modeling the Antennae Galaxies' "antennae" in particular.
A recent study finds that these interacting galaxies are less remote from the Milky Way than previously thought—at 45 million light-years instead of 65 million light-years.
They are located 0.25° north of 31 Crateris and 3.25° southwest of Gamma Corvi.
The Antennae galaxies also contain a relatively young collection of massive globular clusters that were possibly formed as a result of the collision between the two galaxies. The young age of these clusters is in contrast to the average age of most known globular clusters (which are around 12 billion years old), with the formation of the globulars likely originating from shockwaves, generated by the collision of the galaxies, compressing large, massive molecular clouds. The densest regions of the collapsing and compressing clouds are believed to be the birthplace of the clusters.
NGC 4038 has a lot of Cepheid Variables, around 53 of them.
Supernovae
Four supernovae have been observed in NGC 4038.
SN 1974E (type unknown, mag. 14) was discovered by Miklós Lovas on 21 March 1974.
SN 2004gt (type Ic, mag. 14.9) was discovered by Berto Monard on 12 December 2004.
SN 2007sr (type Ia, mag. 12.9) was discovered by the Catalina Sky Survey on 18 December 2007.
SN 2013dk (type Ic, mag. 15.8) was discovered by the CHASE Project (CHilean Automatic Supernovas sEarch) on 22 June 2013.
One supernova has been found in NGC 4039.
SN 1921A (type unknown, mag. 16) was discovered by Edwin Hubble and John Charles Duncan in March 1921.
Timeline
About 1.2 billion years ago, the Antennae were two separate galaxies. NGC 4038 was a barred spiral galaxy and NGC 4039 was a spiral galaxy. 900 million years ago, the Antennae began to approach one another, looking similar to NGC 2207 and IC 2163. 600 million years ago, the Antennae passed through each other, looking like the Mice Galaxies. 300 million years ago, the Antennae's stars began to be released from both galaxies. Today the two streamers of ejected stars extend far beyond the original galaxies, resulting in the antennae shape.
Within 400 million years, the Antennae's nuclei will collide and become a single core with stars, gas, and dust around it. Observations and simulations of colliding galaxies (e.g., by Alar Toomre) suggest that the Antennae Galaxies will eventually form an elliptical galaxy.
X-ray source
Areas containing large amounts of neon (Ne), magnesium (Mg), and silicon (Si) were found when the Chandra X-ray Observatory analyzed the Antennae Galaxies. Heavy elements such as these are necessary in order for planets that may contain life (as we know it) to form. The clouds imaged contain 16 times as much magnesium and 24 times as much silicon as the Sun.
| Physical sciences | Notable galaxies | Astronomy |
1112059 | https://en.wikipedia.org/wiki/Tick-borne%20disease | Tick-borne disease | Tick-borne diseases, which afflict humans and other animals, are caused by infectious agents transmitted by tick bites. They are caused by infection with a variety of pathogens, including rickettsia and other types of bacteria, viruses, and protozoa.
The economic impact of tick-borne diseases is considered to be substantial in humans, and tick-borne diseases are estimated to affect ~80 % of cattle worldwide. Most of these pathogens require passage through vertebrate hosts as part of their life cycle. Tick-borne infections in humans, farm animals, and companion animals are primarily associated with wildlife animal reservoirs. Many tick-borne infections in humans involve a complex cycle between wildlife animal reservoirs and tick vectors. The survival and transmission of these tick-borne viruses are closely linked to their interactions with tick vectors and host cells. These viruses are classified into different families, including Asfarviridae, Reoviridae, Rhabdoviridae, Orthomyxoviridae, Bunyaviridae, and Flaviviridae.
The occurrence of ticks and tick-borne illnesses in humans is increasing. Tick populations are spreading into new areas, in part due to climate change. Tick populations are also affected by changes in the populations of their hosts (e.g. deer, cattle, mice, lizards) and those hosts' predators (e.g. foxes). Diversity and availability of hosts and predators can be affected by deforestation and habitat fragmentation.
Because individual ticks can harbor more than one disease-causing agent, patients can be infected with more than one pathogen at the same time, compounding the difficulty in diagnosis and treatment. As the incidence of tick-borne illnesses increases and the geographic areas in which they are found expand, health workers increasingly must be able to distinguish the diverse, and often overlapping, clinical presentations of these diseases.
18 tick-borne pathogens have been identified in the United States according to the Centers for Disease Control and at least 27 are known globally. New tick-borne diseases have been discovered in the 21st century, due in part to the use of molecular assays and next-generation sequencing.
Prevention
Exposure
Ticks tend to be more active during warmer months, though this varies by geographic region and climate. Areas with woods, bushes, high grass, or leaf litter are likely to have more ticks. Those bitten commonly experience symptoms such as body aches, fever, fatigue, joint pain, or rashes. People can limit their exposure to tick bites by wearing light-colored clothing (including pants and long sleeves), using insect repellent with 20%–30% N,N-Diethyl-3-methylbenzamide (DEET), tucking their pants legs into their socks, checking for ticks frequently, and washing and drying their clothing in a hot dryer.
According to the World Health Organization, tick-to-animal transmission is difficult to prevent because animals do not show visible symptoms; the only effective prevention relies on killing ticks on the livestock production facility.
Symptoms
Ticks also have the potential to induce a motor illness characterized by acute, ascending flaccid paralysis. This condition can be fatal if not treated promptly, affecting both humans and animals. It is mainly associated with certain species of ticks. Symptoms typically ranges from fatigue, numbness in the legs, muscle aches, and, to in some cases, paralysis and other severe neurological manifestations.
Tick-borne diseases (TBD) are a major health threat in the US. The number of pathogens and the burden of disease have been increasing over the last couple decades. With improved diagnostics and surveillance, new pathogens are regularly identified, bettering our understanding of TBDs. Unfortunately, diagnosis of these illnesses remains a challenge, with many TBDs presenting with similar nonspecific symptoms and diagnosis requiring a battery of assays to assess patients adequately. New advanced molecular diagnostic methods, including next-generation sequencing and metagenomics analysis, promise improved detection of novel and emerging pathogens with the ability to detect a litany of potential pathogens with a single assay.
Tick removal
Ticks should be removed as soon as safely possible once discovered. They can be removed either by grasping tweezers as close to the mouth as possible and pulling without rotation; some companies market grooved tools that rotate the hypostome to facilitate removal. Chemical methods to make the tick self-detach, or trying to pull the tick out with one's fingers, are not efficient methods. In Australia and New Zealand, where tick-borne infections are less common than tick reactions, the Australasian Society of Clinical Immunology and Allergy recommends seeking medical assistance or killing ticks in-situ by freezing and then leaving them to fall out to prevent allergic/anaphylactic reactions.
Diagnosis
Diagnosing tick-borne diseases involves a dual approach. Some diagnoses rely on clinical observations and symptom analysis, while others are confirmed through laboratory tests. ticks can transmit a wide range of viruses, many of which are arboviruses. In general, specific laboratory tests are not available for rapid diagnosis of tick-borne diseases. Due to their seriousness, antibiotic treatment is often justified based on clinical presentation alone.
Diagnosing Lyme borreliosis relies on clinical criteria, with a history of a tick bite and associated symptoms being crucial. Laboratory diagnosis follows a 'two-tiered diagnostic protocol,' involving detecting specific antibodies using methods such as immunoenzymatic assays and Western blot tests, preferably with recombinant antigens. While ELISA and Western blot have similar sensitivity, Western blot is more specific due to the identification of specific immunoreactive bands. Seroconversion typically occurs around two weeks after symptom onset, but false positive ELISA results can be linked to poorly reactive antibodies against specific antigens, especially in patients with other infectious and non-infectious diseases.
Tick-borne encephalitis (TBE) presents non-specific clinical features, making laboratory diagnosis crucial. The diagnostic process typically involves identifying specific IgM- and IgG-serum antibodies through enzyme-linked immunosorbent assay (ELISA) since these antibodies are detectable in most cases upon hospitalization.
Treatment
Patients with Lyme disease who are treated with appropriate antibiotics usually recover rapidly and completely. Antibiotics commonly used include doxycycline, amoxicillin, or cefuroxime axetil. For Anaplasmosis, ehrlichiosis and Rocky Mountain spotted fever, Doxycycline is the first line treatment for adults and children of all ages. For babesiosis, a combination therapy with atovaquone and azithromycin is most commonly recommended for treatment of mild to moderate babesiosis. Treatment is usually continued for 7 to 10 days. A combination regimen of oral clindamycin and quinine has also been proven effective, but the rate of adverse reactions is significantly higher with this combination. For Powassan virus, there are no medications for treating Powassan virus infections. Medications, however, can help to relieve symptoms and prevent complications. People with severe disease are typically treated in a hospital where they may be given intravenous fluids, fever-reducing medications, breathing support, and other therapies as needed.
Assessing risk
For a person or pet to acquire a tick-borne disease requires that the individual gets bitten by a tick and that the tick feeds for a sufficient period of time. The feeding time required to transmit pathogens differs for different ticks and different pathogens. Transmission of the bacterium that causes Lyme disease is well understood to require a substantial feeding period. In general, soft ticks (Argasidae) transmit pathogens within minutes of attachment because they feed more frequently, whereas hard ticks (Ixodidae) take hours or days, but the latter are more common and harder to remove.
For an individual to acquire infection, the feeding tick must also be infected. Not all ticks are infected. In most places in the US, 30-50% of deer ticks will be infected with Borrelia burgdorferi (the agent of Lyme disease). Other pathogens are much more rare. Ticks can be tested for infection using a highly specific and sensitive qPCR procedure. Several commercial labs provide this service to individuals for a fee. The Laboratory of Medical Zoology (LMZ), a nonprofit lab at the University of Massachusetts, provides a comprehensive TickReport for a variety of human pathogens and makes the data available to the public. Those wishing to know the incidence of tick-borne diseases in their town or state can search the LMZ surveillance database.
Examples
Major tick-borne diseases include:
Bacterial
Lyme disease or borreliosis
Organism: Borrelia burgdorferi sensu lato (bacterium)
Vector: at least 15 species of ticks in the genus Ixodes, including deer tick (Ixodes scapularis (=I. dammini), I. pacificus, I. ricinus (Europe), I. persulcatus (Asia))
Endemic to: The Americas and Eurasia
Symptoms: Fever, arthritis, neuroborreliosis, erythema migrans, cranial nerve palsy, carditis, fatigue, and influenza-like illness
Treatment: Antibiotics – amoxicillin in pregnant adults and children, doxycycline in other adults
Relapsing fever (tick-borne relapsing fever, different from Lyme disease due to different Borrelia species and ticks)
Organisms: Borrelia species such as B. hermsii, B. parkeri, B. duttoni, B. miyamotoi
Vector: Ornithodoros species
Regions : Primarily in Africa, Spain, Saudi Arabia, Asia in and certain areas of Canada and the western United States
Symptoms: Relapsing fever typically presents as recurring high fevers, flu-like symptoms, headaches, and muscular pain, with less common symptoms including rigors, joint pain, altered mentation, cough, sore throat, painful urination, and rash
Treatment: Antibiotics are the treatment for relapsing fever, with doxycycline, tetracycline, or erythromycin being the treatment of choice.
Typhus Several diseases caused by Rickettsia bacteria (below)
Rocky Mountain spotted fever
Organism: Rickettsia rickettsii
Vector: Wood tick (Dermacentor variabilis), D. andersoni
Region (US): East, Southwest
Vector: Amblyomma cajennense
Region (Brazil): São Paulo, Rio de Janeiro, Minas Gerais.
Symptoms:Fever, headache, altered mental status, myalgia, and rash
Treatment: Antibiotic therapy, typically consisting of doxycycline or tetracycline
Helvetica spotted fever
Organism: Rickettsia helvetica
Region(R. helvetica): Confirmed common in ticks in Sweden, Switzerland, France, and Laos
Vector/region(s)#1: Ixodes ricinus is the main European vector.
Symptoms: Most often small red spots, other symptoms are fever, muscle pain, headache and respiratory problems
Treatment: Broad-spectrum antibiotic therapy is needed, phenoxymethylpenicillin likely is sufficient.
Human granulocytic anaplasmosis (formerly human granulocytic ehrlichiosis or HGE)
Organism: Anaplasma phagocytophilum (formerly Ehrlichia phagocytophilum or Ehrlichia equi)
Vector: Lone star tick (Amblyomma americanum), I. scapularis
Region (US): South Atlantic, South-central
Bartonella: Bartonella transmission rates to humans via tick bite are not well established but Bartonella is common in ticks. For example: 4.76% of 2100 ticks tested in a study in Germany
Tularemia
Organism: Francisella tularensis, A. americanum
Vector: D. variabilis, D. andersoni
Region (US): Southeast, South-central, West, widespread
Viral
Tick-borne meningoencephalitis
Organism: TBEV (FSME) virus, a flavivirus from family Flaviviridae
Vector: deer tick (Ixodes scapularis), Ixodes ricinus (Europe), Ixodes persulcatus (Russia + Asia))
Endemic to: Europe and northern Asia
Powassan virus/deer tick virus
Organism: Powassan virus (POWV), a flavivirus from family Flaviviridae. Lineage 2 POWV is also known as deer tick virus (DTV)
Vector: Ixodes cookei, Ix. scapularis, Ix. marxi, Ix. spinipalpusm, Dermacentor andersoni, and D. variabilis
Endemic to: North America and eastern Russia
Colorado tick fever
Organism: Colorado tick fever virus (CTF), a coltivirus from the Reoviridae
Vector: Dermacentor andersoni
Region: US (West)
Crimean-Congo hemorrhagic fever
Organism: CCHF virus, a nairovirus, from the Bunyaviridae
Vector: Hyalomma marginatum, Rhipicephalus bursa
Region: Southern part of Asia, Northern Africa, Southern Europe
Severe febrile illness
Organism: Heartland virus, a phlebovirus, from the Bunyaviridae
Vector: Lone star tick (Amblyomma americanum)
Region: Missouri and Tennessee, United States
Severe febrile illness, headaches, coma in 1/3 patients
Organism: tentatively Alongshan virus, jingmenvirus group in the flavivirus family
Vector: tick (likely Ixodes persulcatus, Ixodes ricinus), mosquitoes
Region: Inner Mongolia but potentially more widespread
Protozoan
Babesiosis
Organism: Babesia microti, Theileria equi
Vector: Ixodes scapularis (deer tick), I. pacificus (western black-legged tick)
Region (US): Northeast, West Coast
Cytauxzoonosis
Organism: Cytauxzoon felis
Vector: Amblyomma americanum (Lone star tick)
Region (US): South, Southeast
Toxin
Tick paralysis
Cause: Toxin
Vector (US): Dermacentor andersoni (Rocky Mountain wood tick), D. variabilis (American dog tick or wood tick)
Region (US): D. andersoni: East, D. variabilis: East, West coast
Vector (Australia): Ixodes holocyclus (Australian paralysis tick)
Region (Australia): East
Allergies
Alpha-gal allergy - Alpha-gal syndrome is likely caused by a hypersensitivity reaction to the Alpha-gal (Galactose-alpha-1,3-galactose) sugar molecule introduced by ticks while feeding on a human host. The immune reaction can leave people with an allergy to red meat and other mammalian derived products.
The experimental confirmation and investigation of how tick bites contribute to the development of AGS have been established and examined using a mouse model.
| Biology and health sciences | Concepts | Health |
1113241 | https://en.wikipedia.org/wiki/Stegodon | Stegodon | Stegodon ("roofed tooth" from the Ancient Greek words , , 'to cover', + , , 'tooth' because of the distinctive ridges on the animal's molars) is an extinct genus of proboscidean, related to elephants. It was originally assigned to the family Elephantidae along with modern elephants but is now placed in the extinct family Stegodontidae. Like elephants, Stegodon had teeth with plate-like lophs that are different from those of more primitive proboscideans like gomphotheres and mammutids. Fossils of the genus are known from Africa and across much of Asia, as far southeast as Timor (with a single record in southeast Europe). The oldest fossils of the genus are found in Late Miocene strata in Asia, likely originating from the more archaic Stegolophodon, subsequently migrating into Africa. While the genus became extinct in Africa during the Pliocene, Stegodon persisted in South, Southeast and Eastern Asia into the Late Pleistocene.
Morphology
The skull of Stegodon is relatively tall but short, and similar in many respects to those of living elephants. The lower jaw in comparison to early elephantimorphs and its ancestor Stegolophodon is shortened (brevirostrine), and lacks lower tusks/incisors. The molar teeth are superficially like those of elephants, consisting of parallel lamellae that form ridges but are generally relatively low crowned (brachydont), the numbers of ridges are greater in later species. Members of the genus lack permanent premolars. The tusks are proportionally large, with those of the biggest species being among the largest known tusks in proboscideans, with a particularly large tusk of S. ganesa from the Early Pleistocene of India measured to be long, with an estimated mass of approximately , substantially larger than the largest recorded modern elephant tusk. These tusks have a slight upward curvature, and project forwards and parallel to each other, with the tusks often so close together that they are almost touching, such that the trunk would probably have had to rest on top of the tusks rather than be freely hanging between them as in living elephants.
Size
The Chinese S. zdanskyi is suggested to be the largest species, and is known from an old male (50-plus years old) from the Yellow River that is tall and would have weighed approximately in life. It had a humerus long, a femur long, and a pelvis wide. The Indian S. ganesa is suggested to have a shoulder height of about , and a body mass of around . The Javanese species S. trigonocephalus is suggested to have been around tall, with a body mass of around . S. orientalis was around the size of an Asian elephant (Elephas maximus).
Similar to modern-day elephants, stegodonts were likely good swimmers, allowing them to disperse to remote islands in Indonesia, the Philippines and Japan. Once present on the islands, due to the process of insular dwarfism, as a result of decreased land area and the reduction of predation and competition pressure, they reduced in body size, with the degree of dwarfism varying between islands as the result of local conditions. One of the smallest species, Stegodon sumbaensis from Sumba in Indonesia, is estimated at around 8% of the size of mainland Stegodon species, with a body mass of . Sometimes the same island was colonised multiple times by Stegodon, as in Flores, where the Early Pleistocene strongly dwarfed species Stegodon sondaari, which was tall at the shoulder and weighed about , was replaced by the species Stegodon florensis during the Middle Pleistocene which was initially substantially larger, but progressively reduced in size over time, with the earlier subspecies Stegodon florensis florensis from the Middle Pleistocene estimated to be around 50% the size of mainland Stegodon species with a shoulder height of around and a body mass of around 1.7 tons, while the later Stegodon florensis insularis from the Late Pleistocene is estimated to be around 17% the size of mainland Stegodon species, with a shoulder height of around , and a body mass of about .
During Pliocene-Early Pleistocene (from around 4-1 million years ago), a succession of endemic dwarf species of Stegodon, probably representing a single lineage lived in the Japanese archipelago, probably derived from the mainland Chinese S. zdanskyi. In chronological succession these species are Stegodon miensis (4-3 million years ago) Stegodon protoaurorae (3-2 million years ago) and Stegodon aurorae, (2-1 million years ago) which show a progressive size reduction through time, possibly as a result of reducing land area of the Japanese archipelago. The latest and smallest species S. aurorae is estimated to be 25% the size of its mainland ancestor with a body mass of around . S. aurorae also shows morphological straits associated with dwarfism, like shortened limbs.
Ecology
Like modern elephants, but unlike more primitive proboscideans, Stegodon is thought to have chewed using a proal movement (a forward stroke from the back to the front) of the lower jaws. This jaw movement is thought to have evolved independently in elephants and stegodontids. Stegodon populations from the Late Pliocene of the India (including Stegodon insignis) are suggested to have been variable mixed feeders, while those from the earliest Pleistocene (including Stegodon ganesa) of the same region are suggested to have been nearly pure grazers based on isotopic analysis. Based on dental microwear analysis, populations of Stegodon from the Pleistocene of China (Stegodon orientalis and Stegodon huananensis) and mainland southeast Asia (S. orientalis) were found to be browsers, with clear niche differentiation from sympatric Elephas populations, which tended towards mixed feeding (both browsing and grazing), though isotopic analysis of Stegodon cf. orientalis specimens from the late Middle Pleistocene of Thailand suggests that these individuals were mixed feeders that consumed a significant amount of C4 grass. Specimens of Stegodon trigonocephalus from the Early-Middle Pleistocene of Java were found to be mixed feeders to grazers, with a diet similar to that of sympatric Elephas hysudrindicus. The dwarf species from Flores, Stegodon sondaari and Stegodon florensis, are suggested to have been mixed feeders and grazers, respectively, based on stable carbon isotopes. Specimens of Stegodon kaiesensis from the Pliocene of East Africa were found to be browsers to mixed feeders, based on mesowear analysis.
Tracks of a group of Stegodon from the Late Pliocene of Japan suggest that like modern elephants, Stegodon lived in social herds.
On Flores, where dwarf Stegodon species were the only large herbivores, they were likely the main prey of the Komodo dragon.
Taxonomy
In the past, stegodonts were believed to be the ancestors of the true elephants and mammoths, but currently they are believed to have no modern descendants. Stegodon is likely derived from Stegolophodon, an extinct genus known from the Miocene of Asia, with transitional fossils between the two genera known from the Late Miocene of Southeast Asia and Yunnan in South China. Stegodon is more closely related to elephants and mammoths than to mastodons. Like elephants, stegodontids are believed to have derived from gomphotheres.
Phylogeny
The following cladogram shows the placement of the genus Stegodon among other proboscideans, based on hyoid characteristics:
List of species
Stegodon kaisensis Late Miocene – Pliocene, Africa
Stegodon zdanskyi Late Miocene – Pliocene, China
Stegodon huananensis Early Pleistocene, China
Stegodon orientalis Middle – Late Pleistocene, China, Southeast Asia, Japan, Taiwan
Stegodon namadicus/S. insignis/S. ganesa Pliocene – Late Pleistocene, India
Stegodon miensis Pliocene, Japan
Stegodon protoaurorae Late Pliocene – Early Pleistocene, Japan
Stegodon aurorae Early Pleistocene – early Middle Pleistocene, Japan
Stegodon sondaari Early Pleistocene, Flores, Indonesia
Stegodon florensis Middle – Late Pleistocene, Flores, Indonesia
Stegodon luzonensis Middle Pleistocene, Luzon, Philippines
Stegodon trigonocephalus late Early Pleistocene – early Late Pleistocene, Java, Indonesia
Stegodon sompoensis Late Pliocene – Early Pleistocene, Sulawesi, Indonesia
Stegodon sumbaensis Middle – Late Pleistocene, Sumba, Indonesia
Stegodon timorensis Middle Pleistocene, Timor, Indonesia
Stegodon mindanensis Pleistocene Mindanao, Philippines
An indeterminate Stegodon molar of an uncertain locality and age is known from Greece, representing the only record of the genus in Europe. Indeterminate remains are also known from the Early Pleistocene and early Middle Pleistocene of Israel.
Relationship with humans
Remains at a number of sites suggest that humans (in a broad sense, including archaic humans) interacted with Stegodon. At a cave deposit on Gele Mountain near Chongqing in southwest China, a mandible of Stegodon orientalis was used to make a handaxe, with dating suggesting the bone is around 170,000 years old. At the late Middle Pleistocene Panxian Dadong cave site in southern Guizhou Province, southwest China, dating to around 300-190,000 years ago, numerous remains of juvenile (0-12 years of age) and a much smaller number of adult remains of adult Stegodon orientalis, representing a minimum of 12 individuals were found at the site in association with stone tools and human remains. It suggested that Stegodon remains were brought to the cave by humans though none of the elements show clear evidence of processing. At the Xinlong Cave site in the Three Gorges area of Chongqing, suggested to date to around 200-130,000 years ago, two Stegodon cf. orientalis tusks have been found along with human remains. These tusks appear to have been delibrately engraved with patterns and are suggested to have been brought into the cave by humans. At the Late Pleistocene Ma’anshan site also in Guizhou, remains of Stegodon orientalis including both adults and juveniles among other animals are found in two layers, the older dating to around 53,000 years Before Present (BP), with the younger dates to around 19,295-31,155 years BP with the minimum number of individuals being 7 and 2 for the older and younger layers respectively, with the older layer containing adults and juveniles while in the younger later only juveniles are present. Bones at the site display cut marks indicating butchery, and are thought to have been accumulated at the site by people, likely by hunting or possibly scavenging in the case of the large adults found in the older layer.
At Liang Bua cave on Flores dating to around 80-50,000 years ago, remains of the dwarf Stegodon species Stegodon florensis are associated with stone tools produced by the dwarf archaic human species Homo floresiensis, with a small number of the bones bearing cut marks. The ambiguous circumstantial association between bones and stone tools, and the rarity of cut marks makes it unclear to what if to any degree, hunting of Stegodon was practiced by Homo floresiensis.
Evolution and extinction
The oldest fossils of Stegodon in Asia date to the Late Miocene, around 8-11 million years ago, with the oldest fossils of the genus in Africa being around 7-6 million years old. Stegodon became extinct in Africa during the late Pliocene, around 3 million years ago suggested to be the result of expansion of grassland habitats.
The Javanese species Stegodon trigonocephalus became extinct around 130-80,000 years ago during the latest Middle Pleistocene-early Late Pleistocene (Marine Isotope Stage 5) following a change to more humid conditions, which may have reduced grazing habitat. The last records of Stegodon florensis date to around 50,000 years ago, around the time of arrival of modern humans to Flores (the earliest evidence of which dates to 46,000 years ago), suggesting that effects of modern human activity were likely the cause of their extinction.
Stegodon became extinct in the Indian subcontinent (Stegodon namadicus/Stegodon sp.), mainland Southeast Asia and China (S. orientalis) at some point during the Late Pleistocene epoch, while Asian elephants, which existed in sympatry with Stegodon in these regions, are still extant. The precise timing of extinction is uncertain for these regions, though in India records of Stegodon may date as recently as 35-30,000 years ago. The survival of the Asian elephant as opposed to Stegodon orientalis in Southeast Asia and South China has been suggested to be due to its more flexible diet in comparison to S. orientalis. Although some authors have claimed a Holocene survival in China for S. orientalis, these claims cannot be substantiated due to loss of specimens and issues regarding dating.
| Biology and health sciences | Proboscidea | Animals |
4399268 | https://en.wikipedia.org/wiki/East%20Asian%20monsoon | East Asian monsoon | The East Asian monsoon is a monsoonal flow that carries moist air from the Indian Ocean and Pacific Ocean to East Asia. It affects approximately one-third of the global population, influencing the climate of Japan, the Korean Peninsula, Taiwan, China, the Philippines and Mainland Southeast Asia but most significantly Vietnam. It is driven by temperature differences between the East Asian continent and the Pacific Ocean. The East Asian monsoon is divided into a warm and wet summer monsoon and a cold and dry winter monsoon. This cold and dry winter monsoon is responsible for the aeolian dust deposition and pedogenesis that resulted in the creation of the Loess Plateau. The monsoon influences weather patterns as far north as Siberia, causing wet summers that contrast with the cold and dry winters caused by the Siberian High, which counterbalances the monsoon's effect on northerly latitudes.
In most years, the monsoonal flow shifts in a very predictable pattern, with winds being southeasterly in late June, bringing significant rainfall to the region, resulting in the East Asian rainy season as the monsoon boundary advances northward during the spring and summer. This leads to a reliable precipitation spike in July and August. However, this pattern occasionally fails, leading to drought and crop failure. In the winter, the winds are northeasterly and the monsoonal precipitation bands move back to the south, and intense precipitation occurs over southern China and Taiwan.
Over Japan and Korea, the monsoon boundary typically takes the form of a quasi-stationary front separating the cooler air mass associated with the Okhotsk High to the north from the hot, humid air mass associated with the subtropical ridge to the south. After the monsoon boundary passes north of a given location, it is not uncommon for daytime temperatures to exceed with dewpoints of or higher. The spring-summer rainy season is referred to as "plum rain" in various languages of East Asia. In Japan the monsoon boundary is referred to as the as it advances northward during the spring, while it is referred to as the shurin when the boundary retreats back southward during the autumn months. The East Asian monsoon is known as meiyu () in China and Taiwan, and jangma (장마) in Korea.
The location and strength of the East Asian monsoon has varied during the Holocene which scientists track using pollen and dust.
| Physical sciences | Seasons | Earth science |
4402647 | https://en.wikipedia.org/wiki/Notopteridae | Notopteridae | The family Notopteridae contains 10 species of osteoglossiform (bony-tongued) fishes, commonly known as featherbacks and knifefishes. These fishes live in freshwater or brackish environments in Africa and West South East and Southeast Asia.
With the denotation of "knifefish", the notopterids should not be confused with Gymnotiformes, the electric knifefishes from South and Central America. Although their manner of swimming is similar and they are superficially similar in appearance, the two groups are not closely related.
A few of the larger species, especially Chitala ornata, are food fish and occasionally aquarium pets. The name is from Greek noton meaning "back" and pteron meaning "fin".
Fossils
The earliest fossil of this family is otolith of Notopteridarum from the Late Cretaceous of India, about 70.6 to 66 million years ago.
Description
Featherbacks have slender, elongated, bodies, giving them a knife-like appearance. The caudal fin is small and fused with the anal fin, which runs most of the length of the body. Where present, the dorsal fin is small and narrow, giving rise to the common name of "featherback". The fish swims by holding its body rigid and rippling the anal fin to propel itself forward or backwards.
Notopterids have specialized swim bladders. The organ extends throughout the body and even into the fins in some cases. Although the swim bladder is not highly vascularised, it can absorb oxygen from air and also functions to produce sound by squeezing air through a narrow passage into the pharynx.
At least some species prepare nests and guard the eggs until they hatch.
Species
The 10 species in four genera are:
Subfamily Xenomystinae Greenwood 1963 (African knifefishes, African featherbacks)
Genus Papyrocranus Greenwood, 1963
Papyrocranus afer (Günther, 1868) (reticulated knifefish)
Papyrocranus congoensis (Nichols & La Monte, 1932)
Genus Xenomystus Günther, 1868
Xenomystus nigri (Günther, 1868) (African brown knifefish)
Subfamily Notopterinae Bleeker 1851 (Asian knifefishes, Asian featherbacks)
Genus Chitala Fowler, 1934
Chitala blanci (d'Aubenton, 1965) (royal knifefish or Indochina featherback)
Chitala borneensis (Bleeker, 1851) (Indonesian featherback)
Chitala chitala (F. Hamilton, 1822) (Indian featherback)
Chitala hypselonotus (Bleeker, 1852)
†Chitala lopis (Bleeker, 1851)
Chitala ornata (J. E. Gray, 1831) (clown featherback or clown knifefish)
Genus Notopterus Lacépède, 1800
Notopterus notopterus (Pallas, 1769) (bronze featherback)
| Biology and health sciences | Osteoglossiformes | Animals |
4403651 | https://en.wikipedia.org/wiki/Sheep%E2%80%93goat%20hybrid | Sheep–goat hybrid | A sheep–goat hybrid (called a geep in popular media or sometimes a shoat) is a hybrid between a sheep and a goat.
While sheep and goats are similar and can be mated, they belong to different genera in the subfamily Caprinae of the family Bovidae. Sheep belong to the genus Ovis and have 54 chromosomes, while goats belong to the genus Capra and have 60 chromosomes. The offspring of a sheep–goat pairing is generally stillborn. Despite widespread shared pasturing of goats and sheep, hybrids are very rare, demonstrating the genetic distance between the two species. They are not to be confused with sheep–goat chimera, which are artificially created by combining the embryos of a goat and a sheep.
Characteristics
There is a long-standing belief in sheep–goat hybrids, which is presumably due to the animals' resemblance to each other. Some primitive varieties of sheep may be misidentified as goats. In Darwinism – An Exposition of the Theory of Natural Selection with Some of Its Applications (1889), Alfred Russel Wallace wrote:
[...] the following statement of Mr. Low: "It has been long known to shepherds, though questioned by naturalists, that the progeny of the cross between the sheep and goat is fertile. Breeds of this mixed race are numerous in the north of Europe." Nothing appears to be known of such hybrids either in Scandinavia or in Italy; but Professor Giglioli of Florence has kindly given me some useful references to works in which they are described. The following extract from his letter is very interesting: "I need not tell you that there being such hybrids is now generally accepted as a fact. Buffon (Supplements, tom. iii. p. 7, 1756) obtained one such hybrid in 1751 and eight in 1752. Sanson (La Culture, vol. vi. p. 372, 1865) mentions a case observed in the Vosges, France. Geoff. St. Hilaire (Hist. Nat. Gén. des reg. org., vol. iii. p. 163) was the first to mention, I believe, that in different parts of South America the ram is more usually crossed with the she-goat than the sheep with the he-goat. The well-known 'pellones' of Chile are produced by the second and third generation of such hybrids (Gay, 'Hist, de Chile,' vol. i. p. 466, Agriculture, 1862). Hybrids bred from goat and sheep are called 'chabin' in French, and 'cabruno' in Spanish. In Chile such hybrids are called 'carneros lanudos'; their breeding inter se appears to be not always successful, and often the original cross has to be recommenced to obtain the proportion of three-eighths of he-goat and five-eighths of sheep, or of three-eighths of ram and five-eighths of she-goat; such being the reputed best hybrids."
Supposedly, most sheep–goat hybrids die as embryos. Hybrid male mammals are often sterile, demonstrating a phenomenon known as Haldane's rule. The Haldane phenomenon may apply even when the parent species have the same number of chromosomes, as in most cat-species hybrids. It sometimes does not apply when the species chromosome number is different, as in wild horse (chromosome number = 66) with domestic horse (chromosome number = 64) hybrids. Hybrid female fertility tends to decrease with increasing divergence in chromosome similarity between parent species. Presumably, this is due to mismatch problems during meiosis and the resulting production of eggs with unbalanced genetic complements. However, a buck–ewe hybrid born in 2014 died of pregnancy related complications in 2018 raising the question if the parent–species combination has an influence on hybrid fertility.
Blood transcriptome analysis of a buck–ewe hybrid and its parents revealed significant deviations from previously described imprinting schemes and a higher contribution from the goat genome to the genes expressed in the hybrid's blood. Due to the common genome, buck–ewe hybrid share 870 common genes with the maternal goat and 368 genes with the paternal sheep.
Alleged and confirmed cases
At the Botswana Ministry of Agriculture in 2000, a male sheep impregnated a female goat resulting in a live male offspring. This hybrid had 57 chromosomes, intermediate between sheep (54) and goats (60) and was intermediate between the two parent species in type. It had a coarse outer coat, a woolly inner coat, long goat-like legs and a heavy sheep-like body. Although infertile, the hybrid had a very active libido, mounting both ewes and does even when they were not in heat. He was castrated when he was 10 months old, as were the other kids and lambs in the herd.
A male sheep impregnated a female goat in New Zealand resulting in a mixed litter of kids and a female sheep–goat hybrid with 57 chromosomes. The hybrid was subsequently shown to be fertile when mated with a ram.
In France, natural mating of a doe with a ram produced a female hybrid carrying 57 chromosomes. This animal backcrossed in the veterinary college of Nantes to a ram delivered a stillborn and a living male offspring with 54 chromosomes.
In March 2014, a buck–ewe hybrid was born on a farm close to Göttingen in Germany. Also in March 2014, a male buck–ewe hybrid was born in Ireland.
There was a reported case of live births of twin geep on a farm in Ireland in 2018.
There was reported case of a live birth of a sheep–goat hybrid on a farm in Tábor in Czech Republic in 2020. Her name was Barunka and she had health complications after being born. Her owners did not know if she was a goat or sheep, since neither goats nor sheep accepted her. Upon further inspection, it was discovered she was a sheep–goat hybrid.
In May 2021, a healthy doe–ram hybrid was born on a farm in Kentucky, USA despite complications during labor. Her status as a hybrid was confirmed by genetic testing: she has a hybrid karyotype of 57, XX.
Sheep–goat chimera
History
A sheep–goat chimera (sometimes called a geep in popular media) is a chimera produced by combining the embryos of a goat and a sheep; the resulting animal has cells of both sheep and goat origin. A sheep–goat chimera should not be confused with a sheep–goat hybrid, which can result when a goat mates with a sheep. The first sheepgoat chimeras were created by researchers at the Institute of Animal Physiology in Cambridge, England by combining sheep embryos with goat embryos. They reported their results in 1984. The successful chimeras were a mosaic of goat and sheep cells and tissues.
Characteristics
In a chimera, each set of cells (germ line) keeps its own species' identity instead of being intermediate in type between the parental species. Because the chimera contains cells from (at least) two genetically different embryos, and each of these arose by fertilization of an egg by a sperm cell, it has (at least) four genetic parents. In contrast, a hybrid has only two. Although the individual cells in interspecies chimeras are entirely of one of the component species, their behaviour is influenced by the environment in which they find themselves. The sheep–goat chimeras have yielded several demonstrations of this. The most obvious was that the woolly areas of their fleece, tufts of goat wool (angora-type) grew intermingled with ordinary sheep wool, even though the goat breed used in the experiments did not exhibit any wool whatsoever.
Sheepgoat chimeras as a general rule may be assumed to be fertile, with the reservations that apply to chimeras generally (which again reflect that the parent embryos may have been of different sex, so that the animal, apart from being a chimera, may also be intersex). But in accordance with the mosaic- (as opposed to hybrid-) nature of the interspecies chimera, any individual sperm or egg cell it produces must be of either the pure sheep or the pure goat variety. Whether in fact viable germ cells of both species may be, or have been, produced within the same individual is unknown.
The term shoat is sometimes used for sheepgoat hybrids and chimeras, but that term more conventionally means a young piglet.
| Biology and health sciences | Hybrids | Animals |
4404236 | https://en.wikipedia.org/wiki/Bornean%20bearded%20pig | Bornean bearded pig | The Bornean bearded pig (Sus barbatus), also known as the Sunda bearded pig or simply bearded pig, is a species in the pig genus, Sus.
It can be recognized by its prominent beard. It also sometimes has tassels on its tail. It is found in Southeast Asia—Sumatra, Borneo, the Malay Peninsula, and various smaller islands like in Sulu archipelago such as Tawi-Tawi, where it inhabits rainforests and mangrove forests. The bearded pig lives in a family. It can reproduce from the age of 18 months, and can be cross-bred with other species in the family Suidae.
Subspecies
The two subspecies of this pig are:
S. b. barbatus (the nominate subspecies)
S. b. oi (western bearded pig)
As traditionally defined, the nominate is from Borneo. The species is widely ranging in Borneo. It is also found in Tawi-Tawi province at the tip of the Sulu Archipelago in the Philippines, although this population possibly has been extirpated, and S. b. oi is from the Malay Peninsula and Sumatra. Genetic evidence suggests this is incorrect, and S. b. oi is better limited to Sumatra, leaving bearded pigs from both Borneo and the Malay Peninsula in the nominate subspecies. Those from Bangka Island appear somewhat intermediate between the two subspecies.
The Palawan bearded pig (Sus ahoenobarbus) has formerly been considered a subspecies of the bearded pig. However, as indicated by its genetic and morphological distinctness, under the phylogenetic species concept (which does not use subspecies) it needs to be elevated to full species status; while the situation is less clear under other species concepts (as not all S. barbatus populations have been restudied in modern times), the presently available information seems to favor full species status for S. ahoenobarbus in any case.
In captivity
The San Diego Zoo was the first zoo in the Western Hemisphere to breed them.
As of January 2016, it was held in the London Zoo, Berlin Zoo, Gladys Porter Zoo, National Zoo of Malaysia (Zoo Negara), Zoo Taiping, Night Safari, and Southwick's Zoo. The animals at Hellabrunn Zoo were euthanized in 2017 because of old age, and there is only one male left at the Berlin Zoo. Three individuals (one castrated male and two females) left at London Zoo and one individual left at Gladys Porter Zoo and the individuals were replaced by red river hogs in Southwick's Zoo as of 2017, which means that the species will likely disappear soon from European and American zoos.
The death of Neo, the last Bornean bearded pig held in Europe, was announced by Berlin zoo via Facebook on 28 February 2024.
| Biology and health sciences | Pigs_2 | Animals |
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