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24,957,519 | https://en.wikipedia.org/wiki/Ionized%20impurity%20scattering | In quantum mechanics, ionized impurity scattering is the scattering of charge carriers by ionization in the lattice. The most primitive models can be conceptually understood as a particle responding to unbalanced local charge that arises near a crystal impurity; similar to an electron encountering an electric field. This effect is the mechanism by which doping decreases mobility.
In the current quantum mechanical picture of conductivity the ease with which electrons traverse a crystal lattice is dependent on the near perfectly regular spacing of ions in that lattice. Only when a lattice contains perfectly regular spacing can the ion-lattice interaction (scattering) lead to almost transparent behavior of the lattice. Impurity atoms in a crystal have an effect similar to thermal vibrations where conductivity has a direct relationship with temperature.
A crystal with impurities is less regular than a pure crystal, and a reduction in electron mean free paths occurs. Impure crystals have lower conductivity than pure crystals with less temperature sensitivity in that lattice.
See also
Lattice scattering
References
External links
Quantum mechanics
Scattering | Ionized impurity scattering | [
"Physics",
"Chemistry",
"Materials_science"
] | 206 | [
"Theoretical physics",
"Scattering stubs",
"Quantum mechanics",
"Scattering",
"Condensed matter physics",
"Particle physics",
"Nuclear physics",
"Quantum physics stubs"
] |
24,957,935 | https://en.wikipedia.org/wiki/Lattice%20scattering | Lattice scattering is the scattering of ions by interaction with atoms in a lattice. This effect can be qualitatively understood as phonons colliding with charge carriers.
In the current quantum mechanical picture of conductivity the ease with which electrons traverse a crystal lattice is dependent on the near perfectly regular spacing of ions in that lattice. Only when a lattice contains perfectly regular spacing can the ion-lattice interaction (scattering) lead to almost transparent behavior of the lattice.
In the quantum understanding, an electron is viewed as a wave traveling through a medium. When the wavelength of the electrons is larger than the crystal spacing, the electrons will propagate freely throughout the metal without collision.
See also
Ionized impurity scattering
Phonon scattering
References
External links
Quantum mechanics
Scattering | Lattice scattering | [
"Physics",
"Chemistry",
"Materials_science"
] | 158 | [
"Theoretical physics",
"Scattering stubs",
"Quantum mechanics",
"Scattering",
"Condensed matter physics",
"Particle physics",
"Nuclear physics",
"Quantum physics stubs"
] |
24,958,527 | https://en.wikipedia.org/wiki/Group%20testing | In statistics and combinatorial mathematics, group testing is any procedure that breaks up the task of identifying certain objects into tests on groups of items, rather than on individual ones. First studied by Robert Dorfman in 1943, group testing is a relatively new field of applied mathematics that can be applied to a wide range of practical applications and is an active area of research today.
A familiar example of group testing involves a string of light bulbs connected in series, where exactly one of the bulbs is known to be broken. The objective is to find the broken bulb using the smallest number of tests (where a test is when some of the bulbs are connected to a power supply). A simple approach is to test each bulb individually. However, when there are a large number of bulbs it would be much more efficient to pool the bulbs into groups. For example, by connecting the first half of the bulbs at once, it can be determined which half the broken bulb is in, ruling out half of the bulbs in just one test.
Schemes for carrying out group testing can be simple or complex and the tests involved at each stage may be different. Schemes in which the tests for the next stage depend on the results of the previous stages are called adaptive procedures, while schemes designed so that all the tests are known beforehand are called non-adaptive procedures. The structure of the scheme of the tests involved in a non-adaptive procedure is known as a pooling design.
Group testing has many applications, including statistics, biology, computer science, medicine, engineering and cyber security. Modern interest in these testing schemes has been rekindled by the Human Genome Project.
Basic description and terms
Unlike many areas of mathematics, the origins of group testing can be traced back to a single report written by a single person: Robert Dorfman. The motivation arose during the Second World War when the United States Public Health Service and the Selective service embarked upon a large-scale project to weed out all syphilitic men called up for induction. Testing an individual for syphilis involves drawing a blood sample from them and then analysing the sample to determine the presence or absence of syphilis. At the time, performing this test was expensive, and testing every soldier individually would have been very expensive and inefficient.
Supposing there are soldiers, this method of testing leads to separate tests. If a large proportion of the people are infected then this method would be reasonable. However, in the more likely case that only a very small proportion of the men are infected, a much more efficient testing scheme can be achieved. The feasibility of a more effective testing scheme hinges on the following property: the soldiers can be pooled into groups, and in each group the blood samples can be combined. The combined sample can then be tested to check if at least one soldier in the group has syphilis. This is the central idea behind group testing. If one or more of the soldiers in this group has syphilis, then a test is wasted (more tests need to be performed to find which soldier(s) it was). On the other hand, if no one in the pool has syphilis then many tests are saved, since every soldier in that group can be eliminated with just one test.
The items that cause a group to test positive are generally called defective items (these are the broken lightbulbs, syphilitic men, etc.). Often, the total number of items is denoted as and represents the number of defectives if it is assumed to be known.
Classification of group-testing problems
There are two independent classifications for group-testing problems; every group-testing problem is either adaptive or non-adaptive, and either probabilistic or combinatorial.
In probabilistic models, the defective items are assumed to follow some probability distribution and the aim is to minimise the expected number of tests needed to identify the defectiveness of every item. On the other hand, with combinatorial group testing, the goal is to minimise the number of tests needed in a 'worst-case scenario' – that is, create a minmax algorithm – and no knowledge of the distribution of defectives is assumed.
The other classification, adaptivity, concerns what information can be used when choosing which items to group into a test. In general, the choice of which items to test can depend on the results of previous tests, as in the above lightbulb problem. An algorithm that proceeds by performing a test, and then using the result (and all past results) to decide which next test to perform, is called adaptive. Conversely, in non-adaptive algorithms, all tests are decided in advance. This idea can be generalised to multistage algorithms, where tests are divided into stages, and every test in the next stage must be decided in advance, with only the knowledge of the results of tests in previous stages.
Although adaptive algorithms offer much more freedom in design, it is known that adaptive group-testing algorithms do not improve upon non-adaptive ones by more than a constant factor in the number of tests required to identify the set of defective items. In addition to this, non-adaptive methods are often useful in practice because one can proceed with successive tests without first analysing the results of all previous tests, allowing for the effective distribution of the testing process.
Variations and extensions
There are many ways to extend the problem of group testing. One of the most important is called noisy group testing, and deals with a big assumption of the original problem: that testing is error-free. A group-testing problem is called noisy when there is some chance that the result of a group test is erroneous (e.g. comes out positive when the test contained no defectives). The Bernoulli noise model assumes this probability is some constant, , but in general it can depend on the true number of defectives in the test and the number of items tested. For example, the effect of dilution can be modelled by saying a positive result is more likely when there are more defectives (or more defectives as a fraction of the number tested), present in the test. A noisy algorithm will always have a non-zero probability of making an error (that is, mislabeling an item).
Group testing can be extended by considering scenarios in which there are more than two possible outcomes of a test. For example, a test may have the outcomes and , corresponding to there being no defectives, a single defective, or an unknown number of defectives larger than one. More generally, it is possible to consider the outcome-set of a test to be for some .
Another extension is to consider geometric restrictions on which sets can be tested. The above lightbulb problem is an example of this kind of restriction: only bulbs that appear consecutively can be tested. Similarly, the items may be arranged in a circle, or in general, a net, where the tests are available paths on the graph. Another kind of geometric restriction would be on the maximum number of items that can be tested in a group, or the group sizes might have to be even and so on. In a similar way, it may be useful to consider the restriction that any given item can only appear in a certain number of tests.
There are endless ways to continue remixing the basic formula of group testing. The following elaborations will give an idea of some of the more exotic variants. In the 'good–mediocre–bad' model, each item is one of 'good', 'mediocre' or 'bad', and the result of a test is the type of the 'worst' item in the group. In threshold group testing, the result of a test is positive if the number of defective items in the group is greater than some threshold value or proportion. Group testing with inhibitors is a variant with applications in molecular biology. Here, there is a third class of items called inhibitors, and the result of a test is positive if it contains at least one defective and no inhibitors.
History and development
Invention and initial progress
The concept of group testing was first introduced by Robert Dorfman in 1943 in a short report published in the Notes section of Annals of Mathematical Statistics. Dorfman's report – as with all the early work on group testing – focused on the probabilistic problem, and aimed to use the novel idea of group testing to reduce the expected number of tests needed to weed out all syphilitic men in a given pool of soldiers. The method was simple: put the soldiers into groups of a given size, and use individual testing (testing items in groups of size one) on the positive groups to find which were infected. Dorfman tabulated the optimum group sizes for this strategy against the prevalence rate of defectiveness in the population.
Stephen Samuels
found a closed-form solution for the optimal group size as a function of the prevalence rate.
After 1943, group testing remained largely untouched for a number of years. Then in 1957, Sterrett produced an improvement on Dorfman's procedure. This newer process starts by again performing individual testing on the positive groups, but stopping as soon as a defective is identified. Then, the remaining items in the group are tested together, since it is very likely that none of them are defective.
The first thorough treatment of group testing was given by Sobel and Groll in their formative 1959 paper on the subject. They described five new procedures – in addition to generalisations for when the prevalence rate is unknown – and for the optimal one, they provided an explicit formula for the expected number of tests it would use. The paper also made the connection between group testing and information theory for the first time, as well as discussing several generalisations of the group-testing problem and providing some new applications of the theory.
The fundamental result by Peter Ungar in 1960 shows that if the prevalence rate , then individual testing is the optimal group testing procedure with respect to the expected number of tests, and if , then it is not optimal. However, it is important to note that despite 80 years' worth of research effort, the optimal procedure is yet unknown for and a general population size .
Combinatorial group testing
Group testing was first studied in the combinatorial context by Li in 1962, with the introduction of Li’s -stage algorithm. Li proposed an extension of
Dorfman's '2-stage algorithm' to an arbitrary number of stages that required no more
than tests to be guaranteed to find or fewer defectives among items.
The idea was to remove all the items in negative tests, and divide the remaining items into groups as was done with the initial pool. This was to be done times before performing individual testing.
Combinatorial group testing in general was later studied more fully by Katona in 1973. Katona introduced the matrix representation of non-adaptive group-testing and produced a procedure for finding the defective in the non-adaptive 1-defective case in no more than tests, which he also proved to be optimal.
In general, finding optimal algorithms for adaptive combinatorial group testing is difficult, and although the computational complexity of group testing has not been determined, it is suspected to be hard in some complexity class. However, an important breakthrough occurred in 1972, with the introduction of the generalised binary-splitting algorithm. The generalised binary-splitting algorithm works by performing a binary search on groups that test positive, and is a simple algorithm that finds a single defective in no more than the information-lower-bound number of tests.
In scenarios where there are two or more defectives, the generalised binary-splitting algorithm still produces near-optimal results, requiring at most tests above the information lower bound where is the number of defectives. Considerable improvements to this were made in 2013 by Allemann, getting the required number of tests to less than above the information lower bound when and . This was achieved by changing the binary search in the binary-splitting algorithm to a complex set of sub-algorithms with overlapping test groups. As such, the problem of adaptive combinatorial group testing – with a known number or upper bound on the number of defectives – has essentially been solved, with little room for further improvement.
There is an open question as to when individual testing is minmax. Hu, Hwang and Wang showed in 1981 that individual testing is minmax when , and that it is not minmax when . It is currently conjectured that this bound is sharp: that is, individual testing is minmax if and only if . Some progress was made in 2000 by Riccio and Colbourn, who showed that for large , individual testing is minmax when .
Non-adaptive and probabilistic testing
One of the key insights in non-adaptive group testing is that significant gains can be made by eliminating the requirement that the group-testing procedure be certain to succeed (the "combinatorial" problem), but rather permit it to have some low but non-zero probability of mis-labelling each item (the "probabilistic" problem). It is known that as the number of defective items approaches the total number of items, exact combinatorial solutions require significantly more tests than probabilistic solutions — even probabilistic solutions permitting only an asymptotically small probability of error.
In this vein, Chan et al. (2011) introduced COMP, a probabilistic algorithm that requires no more than tests to find up to defectives in items with a probability of error no more than . This is within a constant factor of the lower bound.
Chan et al. (2011) also provided a generalisation of COMP to a simple noisy model, and similarly produced an explicit performance bound, which was again only a constant (dependent on the likelihood of a failed test) above the corresponding lower bound. In general, the number of tests required in the Bernoulli noise case is a constant factor larger than in the noiseless case.
Aldridge, Baldassini and Johnson (2014) produced an extension of the COMP algorithm that added additional post-processing steps. They showed that the performance of this new algorithm, called DD, strictly exceeds that of COMP, and that DD is 'essentially optimal' in scenarios where , by comparing it to a hypothetical algorithm that defines a reasonable optimum. The performance of this hypothetical algorithm suggests that there is room for improvement when , as well as suggesting how much improvement this might be.
Formalisation of combinatorial group testing
This section formally defines the notions and terms relating to group testing.
The input vector, , is defined to be a binary vector of length (that is, ), with the j-th item being called defective if and only if . Further, any non-defective item is called a 'good' item.
is intended to describe the (unknown) set of defective items. The key property of is that it is an implicit input. That is to say, there is no direct knowledge of what the entries of are, other than that which can be inferred via some series of 'tests'. This leads on to the next definition.
Let be an input vector. A set, is called a test. When testing is noiseless, the result of a test is positive when there exists such that , and the result is negative otherwise.
Therefore, the goal of group testing is to come up with a method for choosing a 'short' series of tests that allow to be determined, either exactly or with a high degree of certainty.
A group-testing algorithm is said to make an error if it incorrectly labels an item (that is, labels any defective item as non-defective or vice versa). This is not the same thing as the result of a group test being incorrect. An algorithm is called zero-error if the probability that it makes an error is zero.
denotes the minimum number of tests required to always find defectives among items with zero probability of error by any group-testing algorithm. For the same quantity but with the restriction that the algorithm is non-adaptive, the notation is used.
General bounds
Since it is always possible to resort to individual testing by setting for each , it must be that that . Also, since any non-adaptive testing procedure can be written as an adaptive algorithm by simply performing all the tests without regard to their outcome, . Finally, when , there is at least one item whose defectiveness must be determined (by at least one test), and so .
In summary (when assuming ), .
Information lower bound
A lower bound on the number of tests needed can be described using the notion of sample space, denoted , which is simply the set of possible placements of defectives. For any group testing problem with sample space and any group-testing algorithm, it can be shown that , where is the minimum number of tests required to identify all defectives with a zero probability of error. This is called the information lower bound. This bound is derived from the fact that after each test, is split into two disjoint subsets, each corresponding to one of the two possible outcomes of the test.
However, the information lower bound itself is usually unachievable, even for small problems. This is because the splitting of is not arbitrary, since it must be realisable by some test.
In fact, the information lower bound can be generalised to the case where there is a non-zero probability that the algorithm makes an error. In this form, the theorem gives us an upper bound on the probability of success based on the number of tests. For any group-testing algorithm that performs tests, the probability of success, , satisfies . This can be strengthened to: .
Representation of non-adaptive algorithms
Algorithms for non-adaptive group testing consist of two distinct phases. First, it is decided how many tests to perform and which items to include in each test. In the second phase, often called the decoding step, the results of each group test are analysed to determine which items are likely to be defective. The first phase is usually encoded in a matrix as follows.
Suppose a non-adaptive group testing procedure for items consists of the tests for some . The testing matrix for this scheme is the binary matrix, , where if and only if (and is zero otherwise).
Thus each column of represents an item and each row represents a test, with a in the entry indicating that the test included the item and a indicating otherwise.
As well as the vector (of length ) that describes the unknown defective set, it is common to introduce the result vector, which describes the results of each test.
Let be the number of tests performed by a non-adaptive algorithm. The result vector, , is a binary vector of length (that is, ) such that if and only if the result of the test was positive (i.e. contained at least one defective).
With these definitions, the non-adaptive problem can be reframed as follows: first a testing matrix is chosen, , after which the vector is returned. Then the problem is to analyse to find some estimate for .
In the simplest noisy case, where there is a constant probability, , that a group test will have an erroneous result, one considers a random binary vector, , where each entry has a probability of being , and is otherwise. The vector that is returned is then , with the usual addition on (equivalently this is the element-wise XOR operation). A noisy algorithm must estimate using (that is, without direct knowledge of ).
Bounds for non-adaptive algorithms
The matrix representation makes it possible to prove some bounds on non-adaptive group testing. The approach mirrors that of many deterministic designs, where -separable matrices are considered, as defined below.
A binary matrix, , is called -separable if every Boolean sum (logical OR) of any of its columns is distinct. Additionally, the notation -separable indicates that every sum of any of up to of 's columns is distinct. (This is not the same as being -separable for every .)
When is a testing matrix, the property of being -separable (-separable) is equivalent to being able to distinguish between (up to) defectives. However, it does not guarantee that this will be straightforward. A stronger property, called disjunctness does.
A binary matrix, is called -disjunct if the Boolean sum of any columns does not contain any other column. (In this context, a column A is said to contain a column B if for every index where B has a 1, A also has a 1.)
A useful property of -disjunct testing matrices is that, with up to defectives, every non-defective item will appear in at least one test whose outcome is negative. This means there is a simple procedure for finding the defectives: just remove every item that appears in a negative test.
Using the properties of -separable and -disjunct matrices the following can be shown for the problem of identifying defectives among total items.
The number of tests needed for an asymptotically small average probability of error scales as .
The number of tests needed for an asymptotically small maximum probability of error scales as .
The number of tests needed for a zero probability of error scales as .
Generalised binary-splitting algorithm
The generalised binary-splitting algorithm is an essentially-optimal adaptive group-testing algorithm that finds or fewer defectives among items as follows:
If , test the items individually. Otherwise, set and .
Test a group of size . If the outcome is negative, every item in the group is declared to be non-defective; set and go to step 1. Otherwise, use a binary search to identify one defective and an unspecified number, called , of non-defective items; set and . Go to step 1.
The generalised binary-splitting algorithm requires no more than tests where
.
For large, it can be shown that , which compares favorably to the tests required for Li's -stage algorithm. In fact, the generalised binary-splitting algorithm is close to optimal in the following sense. When it can be shown that , where is the information lower bound.
Non-adaptive algorithms
Non-adaptive group-testing algorithms tend to assume that the number of defectives, or at least a good upper bound on them, is known. This quantity is denoted in this section. If no bounds are known, there are non-adaptive algorithms with low query complexity that can help estimate .
Combinatorial orthogonal matching pursuit (COMP)
Combinatorial Orthogonal Matching Pursuit, or COMP, is a simple non-adaptive group-testing algorithm that forms the basis for the more complicated algorithms that follow in this section.
First, each entry of the testing matrix is chosen i.i.d. to be with probability and otherwise.
The decoding step proceeds column-wise (i.e. by item). If every test in which an item appears is positive, then the item is declared defective; otherwise the item is assumed to be non-defective. Or equivalently, if an item appears in any test whose outcome is negative, the item is declared non-defective; otherwise the item is assumed to be defective. An important property of this algorithm is that it never creates false negatives, though a false positive occurs when all locations with ones in the j-th column of (corresponding to a non-defective item j) are "hidden" by the ones of other columns corresponding to defective items.
The COMP algorithm requires no more than tests to have an error probability less than or equal to . This is within a constant factor of the lower bound for the average probability of error above.
In the noisy case, one relaxes the requirement in the original COMP algorithm that the set of locations of ones in any column of corresponding to a positive item be entirely contained in the set of locations of ones in the result vector. Instead, one allows for a certain number of “mismatches” – this number of mismatches depends on both the number of ones in each column, and also the noise parameter, . This noisy COMP algorithm requires no more than tests to achieve an error probability at most .
Definite defectives (DD)
The definite defectives method (DD) is an extension of the COMP algorithm that attempts to remove any false positives. Performance guarantees for DD have been shown to strictly exceed those of COMP.
The decoding step uses a useful property of the COMP algorithm: that every item that COMP declares non-defective is certainly non-defective (that is, there are no false negatives). It proceeds as follows.
First the COMP algorithm is run, and any non-defectives that it detects are removed. All remaining items are now "possibly defective".
Next the algorithm looks at all the positive tests. If an item appears as the only "possible defective" in a test, then it must be defective, so the algorithm declares it to be defective.
All other items are assumed to be non-defective. The justification for this last step comes from the assumption that the number of defectives is much smaller than the total number of items.
Note that steps 1 and 2 never make a mistake, so the algorithm can only make a mistake if it declares a defective item to be non-defective. Thus the DD algorithm can only create false negatives.
Sequential COMP (SCOMP)
SCOMP (Sequential COMP) is an algorithm that makes use of the fact that DD makes no mistakes until the last step, where it is assumed that the remaining items are non-defective. Let the set of declared defectives be . A positive test is called explained by if it contains at least one item in . The key observation with SCOMP is that the set of defectives found by DD may not explain every positive test, and that every unexplained test must contain a hidden defective.
The algorithm proceeds as follows.
Carry out steps 1 and 2 of the DD algorithm to obtain , an initial estimate for the set of defectives.
If explains every positive test, terminate the algorithm: is the final estimate for the set of defectives.
If there are any unexplained tests, find the "possible defective" that appears in the largest number of unexplained tests, and declare it to be defective (that is, add it to the set ). Go to step 2.
In simulations, SCOMP has been shown to perform close to optimally.
Polynomial Pools (PP)
A deterministic algorithm that is guaranteed to exactly
identify up to positives is Polynomial Pools (PP).
. The algorithm is for the construction of the pooling matrix , which can be straightforwardly used
to decode the observations in . Similar to
COMP, a sample is decoded according to the relation:
, where represents element wise multiplication and is the th column of . Since the decoding step is not difficult, PP is specialized for generating .
Forming groups
A group/pool is
generated using a polynomial relation that specifies the indices of the samples contained in each pool. A set of input parameters determines the algorithm. For a prime number and an integer any prime power is defined by . For a dimension parameter the total number of samples is and the number of samples per pool is .
Further, the
Finite field of order is denoted by
(i.e., the integers defined by special arithmetic operations that ensure that addition and multiplication in remains in ).
The method arranges each sample in a grid and represents it by coordinates . The coordinates are computed according to a polynomial relation using the integers
,
The combination of looping through the values is represented by a set with elements of a sequence of integers, i.e.,
, where
. Without loss of generality, the combination is such that
cycles every times, cycles every times until
cycles only once.
Formulas that compute the sample indices, and thus the corresponding pools, for fixed and , are given by
The computations in can be implemented with publicly available software libraries for finite fields, when is a prime power. When is a prime number then the computations in simplify to modulus arithmetics, i.e., . An example of how to generate one pool when
is displayed in the table below, while the corresponding selection of samples is shown in the figure above.
This method uses tests to exactly identify up to positives among samples.
Because of this PP is particularly effective for large sample sizes, since the number of tests grows only linearly with respect to while the samples grow exponentially with this parameter. However PP can also be effective for small sample sizes.
Example applications
The generality of the theory of group testing lends it to many diverse applications, including clone screening, locating electrical shorts; high speed computer networks; medical examination, quantity searching, statistics; machine learning, DNA sequencing; cryptography; and data forensics. This section provides a brief overview of a small selection of these applications.
Multiaccess channels
A multiaccess channel is a communication channel that connects many users at once. Every user can listen and transmit on the channel, but if more than one user transmits at the same time, the signals collide, and are reduced to unintelligible noise. Multiaccess channels are important for various real-world applications, notably wireless computer networks and phone networks.
A prominent problem with multiaccess channels is how to assign transmission times to the users so that their messages do not collide. A simple method is to give each user their own time slot in which to transmit, requiring slots. (This is called time division multiplexing, or TDM.) However, this is very inefficient, since it will assign transmission slots to users that may not have a message, and it is usually assumed that only a few users will want to transmit at any given time – otherwise a multiaccess channel is not practical in the first place.
In the context of group testing, this problem is usually tackled by dividing time into 'epochs' in the following way. A user is called 'active' if they have a message at the start of an epoch. (If a message is generated during an epoch, the user only becomes active at the start of the next one.) An epoch ends when every active user has successfully transmitted their message. The problem is then to find all the active users in a given epoch, and schedule a time for them to transmit (if they have not already done so successfully). Here, a test on a set of users corresponds to those users attempting a transmission. The results of the test are the number of users that attempted to transmit, and , corresponding respectively to no active users, exactly one active user (message successful) or more than one active user (message collision). Therefore, using an adaptive group testing algorithm with outcomes , it can be determined which users wish to transmit in the epoch. Then, any user that has not yet made a successful transmission can now be assigned a slot to transmit, without wastefully assigning times to inactive users.
Machine learning and compressed sensing
Machine learning is a field of computer science that has many software applications such as DNA classification, fraud detection and targeted advertising. One of the main subfields of machine learning is the 'learning by examples' problem, where the task is to approximate some unknown function when given its value at a number of specific points. As outlined in this section, this function learning problem can be tackled with a group-testing approach.
In a simple version of the problem, there is some unknown function, where , and (using logical arithmetic: addition is logical OR and multiplication is logical AND). Here is ' sparse', which means that at most of its entries are . The aim is to construct an approximation to using point evaluations, where is as small as possible. (Exactly recovering corresponds to zero-error algorithms, whereas is approximated by algorithms that have a non-zero probability of error.)
In this problem, recovering is equivalent to finding . Moreover, if and only if there is some index, , where . Thus this problem is analogous to a group-testing problem with defectives and total items. The entries of are the items, which are defective if they are , specifies a test, and a test is positive if and only if .
In reality, one will often be interested in functions that are more complicated, such as , again where . Compressed sensing, which is closely related to group testing, can be used to solve this problem.
In compressed sensing, the goal is to reconstruct a signal, , by taking a number of measurements. These measurements are modelled as taking the dot product of with a chosen vector. The aim is to use a small number of measurements, though this is typically not possible unless something is assumed about the signal. One such assumption (which is common) is that only a small number of entries of are significant, meaning that they have a large magnitude. Since the measurements are dot products of , the equation holds, where is a matrix that describes the set of measurements that have been chosen and is the set of measurement results. This construction shows that compressed sensing is a kind of 'continuous' group testing.
The primary difficulty in compressed sensing is identifying which entries are significant. Once that is done, there are a variety of methods to estimate the actual values of the entries. This task of identification can be approached with a simple application of group testing. Here a group test produces a complex number: the sum of the entries that are tested. The outcome of a test is called positive if it produces a complex number with a large magnitude, which, given the assumption that the significant entries are sparse, indicates that at least one significant entry is contained in the test.
There are explicit deterministic constructions for this type of combinatorial search algorithm, requiring measurements. However, as with group-testing, these are sub-optimal, and random constructions (such as COMP) can often recover sub-linearly in .
Multiplex assay design for COVID19 testing
During a pandemic such as the COVID-19 outbreak in 2020, virus detection assays are sometimes run using nonadaptive group testing designs.
One example was provided by the Origami Assays project which released open source group testing designs to run on a laboratory standard 96 well plate.
In a laboratory setting, one challenge of group testing is the construction of the mixtures can be time-consuming and difficult to do accurately by hand. Origami assays provided a workaround for this construction problem by providing paper templates to guide the technician on how to allocate patient samples across the test wells.
Using the largest group testing designs (XL3) it was possible to test 1120 patient samples in 94 assay wells. If the true positive rate was low enough, then no additional testing was required.
Data forensics
Data forensics is a field dedicated to finding methods for compiling digital evidence of a crime. Such crimes typically involve an adversary modifying the data, documents or databases of a victim, with examples including the altering of tax records, a virus hiding its presence, or an identity thief modifying personal data.
A common tool in data forensics is the one-way cryptographic hash. This is a function that takes the data, and through a difficult-to-reverse procedure, produces a unique number called a hash. Hashes, which are often much shorter than the data, allow us to check if the data has been changed without having to wastefully store complete copies of the information: the hash for the current data can be compared with a past hash to determine if any changes have occurred. An unfortunate property of this method is that, although it is easy to tell if the data has been modified, there is no way of determining how: that is, it is impossible to recover which part of the data has changed.
One way to get around this limitation is to store more hashes – now of subsets of the data structure – to narrow down where the attack has occurred. However, to find the exact location of the attack with a naive approach, a hash would need to be stored for every datum in the structure, which would defeat the point of the hashes in the first place. (One may as well store a regular copy of the data.) Group testing can be used to dramatically reduce the number of hashes that need to be stored. A test becomes a comparison between the stored and current hashes, which is positive when there is a mismatch. This indicates that at least one edited datum (which is taken as defectiveness in this model) is contained in the group that generated the current hash.
In fact, the amount of hashes needed is so low that they, along with the testing matrix they refer to, can even be stored within the organisational structure of the data itself. This means that as far as memory is concerned the test can be performed 'for free'. (This is true with the exception of a master-key/password that is used to secretly determine the hashing function.)
Notes
References
Citations
General references
Atri Rudra's course on Error Correcting Codes: Combinatorics, Algorithms, and Applications (Spring 2007), Lectures 7.
Atri Rudra's course on Error Correcting Codes: Combinatorics, Algorithms, and Applications (Spring 2010), Lectures 10, 11, 28, 29
See also
Balance puzzle
Combinatorics
Design of experiments | Group testing | [
"Mathematics"
] | 7,649 | [
"Discrete mathematics",
"Combinatorics"
] |
24,958,775 | https://en.wikipedia.org/wiki/P%C3%B6schl%E2%80%93Teller%20potential | In mathematical physics, a Pöschl–Teller potential, named after the physicists Herta Pöschl (credited as G. Pöschl) and Edward Teller, is a special class of potentials for which the one-dimensional Schrödinger equation can be solved in terms of special functions.
Definition
In its symmetric form is explicitly given by
and the solutions of the time-independent Schrödinger equation
with this potential can be found by virtue of the substitution , which yields
.
Thus the solutions are just the Legendre functions with , and , . Moreover, eigenvalues and scattering data can be explicitly computed. In the special case of integer , the potential is reflectionless and such potentials also arise as the N-soliton solutions of the Korteweg–De Vries equation.
The more general form of the potential is given by
Rosen–Morse potential
A related potential is given by introducing an additional term:
See also
Morse potential
Trigonometric Rosen–Morse potential
References list
External links
Eigenstates for Pöschl-Teller Potentials
Quantum mechanical potentials
Mathematical physics
Edward Teller
Quantum models | Pöschl–Teller potential | [
"Physics",
"Mathematics"
] | 229 | [
"Applied mathematics",
"Theoretical physics",
"Quantum mechanics",
"Quantum models",
"Quantum mechanical potentials",
"Mathematical physics",
"Quantum physics stubs"
] |
24,960,160 | https://en.wikipedia.org/wiki/ProVerif | ProVerif is a software tool for automated reasoning about the security properties of cryptographic protocols. The tool has been developed by Bruno Blanchet and others.
Support is provided for cryptographic primitives including: symmetric & asymmetric cryptography; digital signatures; hash functions; bit-commitment; and signature proofs of knowledge. The tool is capable of evaluating reachability properties, correspondence assertions and observational equivalence. These reasoning capabilities are particularly useful to the computer security domain since they permit the analysis of secrecy and authentication properties. Emerging properties such as privacy, traceability and verifiability can also be considered. Protocol analysis is considered with respect to an unbounded number of sessions and an unbounded message space. The tool is capable of attack reconstruction: when a property cannot be proved, an execution trace which falsifies the desired property is constructed.
Applicability of ProVerif
ProVerif has been used in the following case studies, which include the security analysis of actual network protocols:
Abadi & Blanchet used correspondence assertions to verify the certified email protocol.
Abadi, Blanchet & Fournet analyse the Just Fast Keying protocol, which was one of the candidates to replace Internet Key Exchange (IKE) as the key exchange protocol in IPsec, by combining manual proofs with ProVerif proofs of correspondence and equivalence.
Blanchet & Chaudhuri studied the integrity of the Plutus file system on untrusted storage, using correspondence assertions, resulting in the discovery, and subsequent fixing, of weaknesses in the initial system.
Bhargavan et al. use ProVerif to analyse cryptographic protocol implementations written in F#; in particular the Transport Layer Security (TLS) protocol has been studied in this manner.
Chen & Ryan have evaluated authentication protocols found in the Trusted Platform Module (TPM), a widely deployed hardware chip, and discovered vulnerabilities.
Delaune, Kremer & Ryan and Backes, Hritcu & Maffei formalise and analyse privacy properties for electronic voting using observational equivalence.
Delaune, Ryan & Smyth and Backes, Maffei & Unruh analyse the anonymity properties of the trusted computing scheme Direct Anonymous Attestation (DAA) using observational equivalence.
Kusters & Truderung examine protocols with Diffie-Hellman exponentiation and XOR.
Smyth, Ryan, Kremer & Kourjieh formalise and analyse verifiability properties for electronic voting using reachability.
Google verified its transport layer protocol ALTS.
Sardar et al. verified the remote attestation protocols in Intel SGX.
Further examples can be found online: .
Alternatives
Alternative analysis tools include: AVISPA (for reachability and correspondence assertions), KISS (for static equivalence), YAPA (for static equivalence). CryptoVerif for verification of security against polynomial time adversaries in the computational model. The Tamarin Prover is a modern alternative to ProVerif, with excellent support for Diffie-Hellman equational reasoning, and verification of observational equivalence properties.
References
External links
Cryptographic software
Free software programmed in OCaml
Automated reasoning
Software using the GNU General Public License
Software using the BSD license | ProVerif | [
"Mathematics"
] | 672 | [
"Cryptographic software",
"Mathematical software"
] |
24,962,428 | https://en.wikipedia.org/wiki/Rhamnogalacturonan-II | Rhamnogalacturonan-II (RG-II) is a complex polysaccharide component of pectin that is found in the primary cell walls of dicotyledonous and monocotyledonous plants and gymnosperms. It is supposed to be crucial for the plant cell wall integrity. RG-II is also likely to be present in the walls of some lower plants (ferns, horsetails, and lycopods). Its global structure is conserved across vascular plants, albeit a number of variations within the RGII side chains have been observed between different plants. RG-II is composed of 12 different glycosyl residues including D-rhamnose, D-apiose, D-galactose, L-galactose, Kdo, D-galacturonic acid, L-arabinose, D-xylose, and L-aceric acid, linked together by at least 21 distinct glycosidic linkages. Some resides are further modified via methylation and acetylation. It moreover supports borate mediated cross-linking between different RGII side-chain apiosyl residues. The backbone consists of a linear polymer of alpha-1,4-linked D-galactopyranosiduronic acid. RG-II can be isolated from different sources, such as apple juice and red wine.
The gut bacterium Bacteroides thetaiotaomicron has a polysaccharide utilization locus that contains enzymes that allows deconstruction of rhamnogalacturonan-II, cleaving all but 1 of its 21 distinct glycosidic linkages.
See also
Pectin
References
Polysaccharides
Wine chemistry | Rhamnogalacturonan-II | [
"Chemistry"
] | 364 | [
"Carbohydrates",
"Organic compounds",
"Wine chemistry",
"Alcohol chemistry",
"Organic compound stubs",
"Organic chemistry stubs",
"Polysaccharides"
] |
24,969,173 | https://en.wikipedia.org/wiki/Types%20of%20concrete | Concrete is produced in a variety of compositions, finishes and performance characteristics to meet a wide range of needs.
Mix design
Modern concrete mix designs can be complex. The choice of a concrete mix depends on the need of the project both in terms of strength and appearance and in relation to local legislation and building codes.
The design begins by determining the requirements of the concrete. These requirements take into consideration the weather conditions that the concrete will be exposed to in service, and the required design strength. The compressive strength of a concrete is determined by taking standard molded, standard-cured cylinder samples.
Many factors need to be taken into account, from the cost of the various additives and aggregates, to the trade offs between the "slump" for easy mixing and placement and ultimate performance.
A mix is then designed using cement (Portland or other cementitious material), coarse and fine aggregates, water and chemical admixtures. The method of mixing will also be specified, as well as conditions that it may be used in.
This allows a user of the concrete to be confident that the structure will perform properly.
Various types of concrete have been developed for specialist application and have become known by these names.
Concrete mixes can also be designed using software programs. Such software provides the user an opportunity to select their preferred method of mix design and enter the material data to arrive at proper mix designs.
Historic concrete composition
Concrete has been used since ancient times. Regular Roman concrete for example was made from volcanic ash (pozzolana), and hydrated lime. Roman concrete was superior to other concrete recipes (for example, those consisting of only sand and lime) used by other cultures. Besides volcanic ash for making regular Roman concrete, brick dust can also be used. Besides regular Roman concrete, the Romans also invented hydraulic concrete, which they made from volcanic ash and clay.
Some types of concrete used to make garden sculptures and planters have been called composition stone or composite stone. There is no single precise formula that differentiates composition stone from other lime-cemented concretes, which is unsurprising because the term predates modern chemical science, being attested since at latest the 1790s. In the 19th and later centuries, the term artificial stone has encompassed various human-made stones including numerous cemented concretes.
Modern concrete
Regular concrete is the lay term for concrete that is produced by following the mixing instructions that are commonly published on packets of cement, typically using sand or other common material as the aggregate, and often mixed in improvised containers. The ingredients in any particular mix depends on the nature of the application. Regular concrete can typically withstand a pressure from about 10 MPa (1450 psi) to 40 MPa (5800 psi), with lighter duty uses such as blinding concrete having a much lower MPa rating than structural concrete. Many types of pre-mixed concrete are available which include powdered cement mixed with an aggregate, needing only water.
Typically, a batch of concrete can be made by using 1 part Portland cement, 2 parts dry sand, 3 parts dry stone, 1/2 part water. The parts are in terms of weight – not volume. For example, of concrete would be made using cement, water, dry sand, dry stone (1/2" to 3/4" stone). This would make of concrete and would weigh about . The sand should be mortar or brick sand (washed and filtered if possible) and the stone should be washed if possible. Organic materials (leaves, twigs, etc.) should be removed from the sand and stone to ensure the highest strength.
High-strength concrete
High-strength concrete has a compressive strength greater than 40 MPa (6000 psi). In the UK, BS EN 206-1 defines High strength concrete as concrete with a compressive strength class higher than C50/60. High-strength concrete is made by lowering the water-cement (W/C) ratio to 0.35 or lower. Often silica fume is added to prevent the formation of free calcium hydroxide crystals in the cement matrix, which might reduce the strength at the cement-aggregate bond.
Low W/C ratios and the use of silica fume make concrete mixes significantly less workable, which is particularly likely to be a problem in high-strength concrete applications where dense rebar cages are likely to be used. To compensate for the reduced workability, superplasticizers are commonly added to high-strength mixtures. Aggregate must be selected carefully for high-strength mixes, as weaker aggregates may not be strong enough to resist the loads imposed on the concrete and cause failure to start in the aggregate rather than in the matrix or at a void, as normally occurs in regular concrete.
In some applications of high-strength concrete the design criterion is the elastic modulus rather than the ultimate compressive strength.
Stamped concrete
Stamped concrete is an architectural concrete that has a superior surface finish. After a concrete floor has been laid, floor hardeners (can be pigmented) are impregnated on the surface and a mold that may be textured to replicate a stone / brick or even wood is stamped on to give an attractive textured surface finish. After sufficient hardening, the surface is cleaned and generally sealed to provide protection. The wear resistance of stamped concrete is generally excellent and hence found in applications like parking lots, pavements, walkways etc.
High-performance concrete
High-performance concrete (HPC) is a relatively new term for concrete that conforms to a set of standards above those of the most common applications, but not limited to strength. While all high-strength concrete is also high-performance, not all high-performance concrete is high-strength. Some examples of such standards currently used in relation to HPC are:
Ease of placement – HPC can be consolidated adequately by gravity (self consolidating) and fills gaps between bars without vibration.
Compaction without segregation
Early age strength
Long-term mechanical properties
Permeability
Density
Heat of hydration
Toughness
Volume stability
Long life in severe environments
Depending on its implementation, environmental
HPC is concrete that develops a strength greater than at 28, 56, or 90 days. These strengths generally require well-graded hard rock aggregates, a fairly high proportion of cement plus fly ash, water-reducing admixtures, and the silica fume, with relatively low water content. Extended mixing may be necessary to adequately disperse the silica fume, which is generally supplied in a granular format. The rich mixes may cause high heat of hydration in thick placements, which can be moderated by using a higher proportion of fly-ash, up to 30% of the cement content. Limestone powder may also be used to increase fluidity.
Ultra-high-performance concrete
Ultra-high-performance concrete is a new type of concrete that is being developed by agencies concerned with infrastructure protection. UHPC is characterized by being a steel fibre-reinforced cement composite material with compressive strengths in excess of 150 MPa, up to and possibly exceeding 250 MPa.
UHPC is also characterized by its constituent material make-up: typically fine-grained sand, fumed silica, small steel fibers, and special blends of high-strength Portland cement. Note that there is no large aggregate. The current types in production (Ductal, Taktl, etc.) differ from normal concrete in compression by their strain hardening, followed by sudden brittle failure. Ongoing research into UHPC failure via tensile and shear failure is being conducted by multiple government agencies and universities around the world.
Micro-reinforced ultra-high-performance concrete
Micro-reinforced ultra-high-performance concrete is the next generation of UHPC. In addition to high compressive strength, durability and abrasion resistance of UHPC, micro-reinforced UHPC is characterized by extreme ductility, energy absorption and resistance to chemicals, water and temperature. The continuous, multi-layered, three dimensional micro-steel mesh exceeds UHPC in durability, ductility and strength. The performance of the discontinuous and scattered fibers in UHPC is relatively unpredictable. Micro-reinforced UHPC is used in blast, ballistic and earthquake resistant construction, structural and architectural overlays, and complex facades.
Ducon was the early developer of micro-reinforced UHPC, which has been used in the construction of new World Trade Center in New York.
Low-density structural concrete
Ceramic aggregates with a density below that of water are used for low density structural concrete. These aggregates may include expanded clays and shales, preferably with water absorption below 10%. For structural concrete only coarse low density aggregates are used, with natural sand as the fine aggregates. However, lower percentages are used for moderate density concretes.
The concrete can develop high compressive and tensile strengths, while shrinkage and creep remain acceptable, but will generally be less rigid than conventional mixes. The most obvious advantage is the low density, but these concretes also have low permeability to water and greater thermal insulation. Resistance to abrasion by ice is similar to normal concrete. Disadvantages are that the water absorption by the aggregates may be relatively high, and vibrational consolidation can cause the low density aggregate to float. This can be avoided by minimising vibration and using fluid mixes. Low density has advantages for floating structures.
Self-consolidating concrete
The defects in concrete in Japan were found to be mainly due to high water-cement ratio to increase workability. Poor compaction occurred mostly because of the need for speedy construction in the 1960s and 1970s. Hajime Okamura envisioned the need for concrete which is highly workable and does not rely on the mechanical force for compaction. During the 1980s, Okamura and his Ph.D. student Kazamasa Ozawa at the University of Tokyo developed self-compacting concrete (SCC) which was cohesive, but flowable and took the shape of the formwork without use of any mechanical compaction. SCC is known as self-consolidating concrete in the United States.
SCC is characterized by the following:
extreme fluidity as measured by flow, typically between 650–750 mm on a flow table, rather than slump (height)
no need for vibrators to compact the concrete
easier placement
no bleeding, or aggregate segregation
increased liquid head pressure, which can be detrimental to safety and workmanship
SCC can save up to 50% in labor costs due to 80% faster pouring and reduced wear and tear on formwork.
In 2005, self-consolidating concretes accounted for 10–15% of concrete sales in some European countries. In the precast concrete industry in the U.S., SCC represents over 75% of concrete production. 38 departments of transportation in the US accept the use of SCC for road and bridge projects.
This emerging technology is made possible by the use of polycarboxylates plasticizer instead of older naphthalene-based polymers, and viscosity modifiers to address aggregate segregation.
Vacuum concrete
Vacuum concrete, made by using steam to produce a vacuum inside a concrete mixing truck to release air bubbles inside the concrete, is being researched. The idea is that the steam displaces the air normally over the concrete. When the steam condenses into water it will create a low pressure over the concrete that will pull air from the concrete. This will make the concrete stronger due to there being less air in the mixture. A drawback is that the mixing has to be done in an airtight container.
The final strength of concrete is increased by about 25%. Vacuum concrete stiffens very rapidly so that the formworks can be removed within 30 minutes of casting even on columns of 20 ft. high. This is of considerable economic value, particularly in a precast factory as the forms can be reused at frequent intervals. The bond strength of vacuum concrete is about 20% higher. The surface of vacuum concrete is entirely free from pitting and the uppermost 1/16 inch is highly resistant to abrasion. These characteristics are of special importance in the construction of concrete structures which are to be in contact with flowing water at a high velocity.
It bonds well to old concrete and can, therefore, be used for resurfacing road slabs and other repair work.
Shotcrete
Shotcrete (also known by the trade name Gunite) uses compressed air to shoot concrete onto (or into) a frame or structure. The greatest advantage of the process is that shotcrete can be applied overhead or on vertical surfaces without formwork. It is often used for concrete repairs or placement on bridges, dams, pools, and on other applications where forming is costly or material handling and installation is difficult. Shotcrete is frequently used against vertical soil or rock surfaces, as it eliminates the need for formwork. It is sometimes used for rock support, especially in tunneling. Shotcrete is also used for applications where seepage is an issue to limit the amount of water entering a construction site due to a high water table or other subterranean sources. This type of concrete is often used as a quick fix for weathering for loose soil types in construction zones.
There are two application methods for shotcrete.
dry-mix – the dry mixture of cement and aggregates is filled into the machine and conveyed with compressed air through the hoses. The water needed for the hydration is added at the nozzle.
wet-mix – the mixes are prepared with all necessary water for hydration. The mixes are pumped through the hoses. At the nozzle compressed air is added for spraying.
For both methods additives such as accelerators and fiber reinforcement may be used.
Limecrete
In limecrete, lime concrete or roman concrete the cement is replaced by lime.
One successful formula was developed in the mid-1800s by Dr. John E. Park.
Lime has been used since Roman times either as mass foundation concretes or as lightweight concretes using a variety of aggregates combined with a wide range of pozzolans (fired materials) that help to achieve increased strength and speed of set. Lime concrete was used to build monumental architecture during and after the roman concrete revolution as well as a wide variety of applications such as floors, vaults or domes.
Over the last decade, there has been a renewed interest in using lime for these applications again.
Environmental Benefits
Lime is burnt at a lower temperature than cement and so has an immediate energy saving of 20% (although kilns etc. are improving so figures do change). A standard lime mortar has about 60-70% of the embodied energy of a cement mortar. It is also considered to be more environmentally friendly because of its ability, through carbonation, to re-absorb its own weight in Carbon Dioxide (compensating for that given off during burning).
Lime mortars allow other building components such as stone, wood and bricks to be reused and recycled because they can be easily cleaned of mortar and limewash.
Lime enables other natural and sustainable products such as wood (including woodfibre, wood wool boards), hemp, straw etc. to be used because of its ability to control moisture (if cement were used, these buildings would compost).
Health Benefits
Lime plaster is hygroscopic (literally means 'water seeking') which draws the moisture from the internal to the external environment, this helps to regulate humidity creating a more comfortable living environment as well as helping to control condensation and mould growth which have been shown to have links to allergies and asthmas.
Lime plasters and limewash are non-toxic, therefore they do not contribute to indoor air pollution unlike some modern paints.
Pervious concrete
Pervious concrete, used in permeable paving, contains a network of holes or voids, to allow air or water to move through the concrete
This allows water to drain naturally through it, and can both remove the normal surface-water drainage infrastructure, and allow replenishment of groundwater when conventional concrete does not.
It is formed by leaving out some or all of the fine aggregate (fines). The remaining large aggregate then is bound by a relatively small amount of Portland cement. When set, typically between 15% and 25% of the concrete volume is voids, allowing water to drain at around 5 gal/ft2/ min (70 L/m2/min) through the concrete.
Installation
Pervious concrete is installed by being poured into forms, then screeded off, to level (not smooth) the surface, then packed or tamped into place. Due to the low water content and air permeability, within 5–15 minutes of tamping, the concrete must be covered with a 6-mil poly plastic, or it will dry out prematurely and not properly hydrate and cure.
Characteristics
Pervious concrete can significantly reduce noise, by allowing air to be squeezed between vehicle tyres and the roadway to escape. This product cannot be used on major U.S. state highways currently due to the high psi ratings required by most states. Pervious concrete has been tested up to 4500 psi so far.
Cellular concrete
Aerated concrete produced by the addition of an air-entraining agent to the concrete (or a lightweight aggregate such as expanded clay aggregate or cork granules and vermiculite) is sometimes called cellular concrete, lightweight aerated concrete, variable density concrete, Foam Concrete and lightweight or ultra-lightweight concrete, not to be confused with aerated autoclaved concrete, which is manufactured off-site using an entirely different method.
In the 1977 work A Pattern Language: Towns, Buildings and Construction, architect Christopher Alexander wrote in pattern 209 on "Good Materials":
The variable density is normally described in kg per m3, where regular concrete is 2400 kg/m3. Variable density can be as low as 300 kg/m3, although at this density it would have no structural integrity at all and would function as a filler or insulation use only. The variable density reduces strength to increase thermal and acoustical insulation by replacing the dense heavy concrete with air or a light material such as clay, cork granules and vermiculite. There are many competing products that use a foaming agent that resembles shaving cream to mix air bubbles in with the concrete. All accomplish the same outcome: to displace concrete with air.
Applications of foamed concrete include:
Roof insulation
Blocks and panels for walls
Levelling floors
Void filling
Road sub-bases and maintenance
Bridge abutments and repairs
Ground stabilisation
Cork-cement composites
Waste Cork granules are obtained during production of bottle stoppers from the treated bark of Cork oak. These granules have a density of about 300 kg/m3, lower than most lightweight aggregates used for making lightweight concrete. Cork granules do not significantly influence cement hydration, but cork dust may. Cork cement composites have several advantages over standard concrete, such as lower thermal conductivities, lower densities and good energy absorption characteristics. These composites can be made of density from 400 to 1500 kg/m3, compressive strength from 1 to 26 MPa, and flexural strength from 0.5 to 4.0 MPa.
Roller-compacted concrete
Roller-compacted concrete, sometimes called rollcrete, is a low-cement-content stiff concrete placed using techniques borrowed from earthmoving and paving work. The concrete is placed on the surface to be covered, and is compacted in place using large heavy rollers typically used in earthwork. The concrete mix achieves a high density and cures over time into a strong monolithic block. Roller-compacted concrete is typically used for concrete pavement, but has also been used to build concrete dams, as the low cement content causes less heat to be generated while curing than typical for conventionally placed massive concrete pours.
Glass concrete
The use of recycled glass as aggregate in concrete has become popular in modern times, with large scale research being carried out at Columbia University in New York. This greatly enhances the aesthetic appeal of the concrete. Recent research findings have shown that concrete made with recycled glass aggregates have shown better long-term strength and better thermal insulation due to its better thermal properties of the glass aggregates.
Asphalt concrete
Strictly speaking, asphalt is a form of concrete as well, with bituminous materials replacing cement as the binder.
Rapid strength concrete
This type of concrete is able to develop high resistance within few hours after being manufactured. This feature has advantages such as removing the formwork early and to move forward in the building process very quickly, repaired road surfaces that become fully operational in just a few hours. Ultimate strength and durability can vary from that of standard concrete, depending on compositional details.
Rubberized concrete
While "rubberized asphalt concrete" is common, rubberized Portland cement concrete ("rubberized PCC") is still undergoing experimental tests, as of 2009.
Nanoconcrete
Nanoconcrete contains Portland cement particles that are no greater than 100 μm. It is a product of high-energy mixing (HEM) of cement, sand and water.
Polymer concrete
Polymer concrete is concrete which uses polymers to bind the aggregate. Polymer concrete can gain a lot of strength in a short amount of time. For example, a polymer mix may reach 5000 psi in only four hours. Polymer concrete is generally more expensive than conventional concretes.
Geopolymer concrete
Geopolymer cement is an alternative to ordinary Portland cement and is used to produce Geopolymer concrete by adding regular aggregates to a geopolymer cement slurry. It is made from inorganic aluminosilicate (Al-Si) polymer compounds that can utilise recycled industrial waste (e.g. fly ash, blast furnace slag) as the manufacturing inputs resulting in up to 80% lower carbon dioxide emissions. Greater chemical and thermal resistance, and better mechanical properties, are said to be achieved for geopolymer concrete at both atmospheric and extreme conditions.
Similar concretes have not only been used in Ancient Rome (see Roman concrete), but also in the former Soviet Union in the 1950s and 1960s. Buildings in Ukraine are still standing after 45 years.
Refractory cement
High-temperature applications, such as masonry ovens and the like, generally require the use of a refractory cement; concretes based on Portland cement can be damaged or destroyed by elevated temperatures, but refractory concretes are better able to withstand such conditions. Materials may include calcium aluminate cements, fire clay, ganister and minerals high in aluminium.
Innovative mixtures
On-going research into alternative mixtures and constituents has identified potential mixtures that promise radically different properties and characteristics.
Bendable, self-healing concrete
Researchers at the University of Michigan have developed Engineered Cement Composites (ECC), a fiber-reinforced bendable concrete. The composite contains many of the ingredients used in regular concrete, but instead of coarse aggregate it includes microscale fibers.
The mixture has much smaller crack propagation that does not suffer the usual cracking and subsequent loss of strength at high levels of tensile stress. Researchers have been able to take mixtures beyond 3 percent strain, past the more typical 0.1% point at which failure occurs.
In addition, the composition of the material supports self-healing. When cracks occur, extra dry cement in the concrete is exposed. It reacts with water and carbon dioxide to form calcium carbonate and fix the crack.
CO2 sequestering concretes
Researchers have tried to sequester CO2 in concrete by developing advanced materials. One approach is to use magnesium silicate (talc) as an alternative to calcium. This lowers the temperature required for the production process and decreases the release of carbon dioxide during firing. During the hardening phase, additional carbon is sequestered.
A related approach is mineral carbonation (MC). It produces stable carbonate aggregates from calcium- or magnesium-containing materials and CO2. Stable aggregates can be used for concrete or to produce carbon neutral building blocks such as bricks or precast concrete. CarbonCure Technologies uses waste CO2 from oil refineries to make its bricks and wet cement mix, offsetting up to 5% of its carbon footprint. Solidia Technologies fires its brick and precast concrete at lower temperatures and cures them with CO2 gas, claiming to reduce its carbon emissions by 30%. Carbonaide uses carbon dioxide in the curing phase of precast concrete production and has demonstrated up to 40% savings in cement consumption with their first client.
Another method of calcium-based mineral carbonation has been inspired by biomimicry of naturally occurring calcium structures. Ginger Krieg Dosier of bioMASON has developed a method for producing bricks without firing kilns or significant carbon release. The bricks are grown in molds over four days through a process of microbiologically induced calcite precipitation. Sporosarcina pasteurii bacteria forms calcite from water, calcium, and urea, incorporating CO2 from the urea, and releasing ammonia for fertilizer.
One research team found a way to use a form of microalgae called coccolithophores to mass produce calcium carbonate via photosynthesis at a faster rate than corals. They can survive in warm, cold , salt and fresh water. The technique has the potential to absorb more CO2 than it emits. Between 1-2 million acres of open ponds could supply enough microalgae to satisfy US cement consumption. The team claims the material can be immediately substituted into existing production processes.
Living walls resisting dessiccation
Another approach involves the development of bioreceptive lightweight concrete which can be used to create living walls resisting dessiccation. Researchers at the Bartlett School of Architecture are developing materials aimed to support the growth of poikilohydric plants such as algae, mosses and lichens (organisms having no mechanism to prevent desiccation). Once established, the combination of new materials and plants can potentially improve storm-water management and absorb pollutants.
Smog eating
Titanium dioxide has been added to concrete mixtures to reduce smog. A daylight photo-catalytic between the titanium in this concrete and the smog reduces bacteria and dirt from accumulating on the surface. It can also be used to break down nitrogen dioxides created by industrial processes.
Gypsum concrete
Gypsum concrete is a building material used as a floor underlay used in wood-frame and concrete construction for fire ratings, sound reduction, radiant heating, and floor leveling. It is a mixture of gypsum, Portland cement, and sand. One of its advantages is the lightweight nature. It weighs less than regular concrete while maintaining comparable compressive strength and costs. It is also easy to work with and level, allowing for faster installation and higher productivity. The use of gypsum concrete for radiant heat flooring became popular in the 1980s with the introduction of plastic PEX tubing, which is not susceptible to corrosion from the concrete.
Foam concrete
Foam concrete, also known as lightweight cellular concrete or foamed cement, is a cement-based material that incorporates stable air bubbles to create a lightweight and highly insulating product. Unlike air-entrained concrete, which introduces tiny air bubbles through an admixture during mixing, foam concrete replaces coarse aggregates with these air bubbles, resulting in a significant difference in density, with foam concrete typically ranging from 400 kg/m3 to 1600 kg/m3, whereas air-entrained concrete maintains its density. Foam concrete is produced by mixing cement or fly ash, sand, water, and a synthetic aerated foam, which provides stability to the air bubbles, in contrast to air-entrained concrete which is produced incorporating specialized admixtures directly into the concrete mix. Foam concrete offers excellent thermal and acoustic insulation properties, making it suitable for applications such as insulation, void filling, and trench reinstatement. Its lightweight nature also makes it easier to handle and transport compared to traditional concrete. Foam concrete can be easily molded into various shapes and sizes, allowing for versatile applications. Its properties make it suitable for insulation, void filling, and other construction applications where weight reduction and thermal insulation are desired.
Air-entrained concrete
Air-entrained concrete is a type of concrete that intentionally incorporates tiny air bubbles (10 to 500 micrometres in diameter) through the addition of an air entraining agent during the mixing process. These air bubbles enhance the workability of the concrete during placement and improve its durability when hardened, particularly in regions prone to freeze-thaw cycles. Unlike foam concrete, which is lightweight and created by introducing stable air bubbles using a foam agent, air-entrained concrete maintains its density (air consists of 6–12 vol.%) while enhancing durability, workability, and resistance to freeze-thaw cycles. The main benefits of air-entrained concrete include improved workability during placement, increased resistance to cracking and surface damage, enhanced durability against fire damage, and overall strength. Additionally, the air voids in air-entrained concrete act as internal cushioning, absorbing energy during impact and increasing resistance to physical forces, thereby increasing its overall durability.
Marine habitat concrete
Marine habitat concrete is concrete used in artificial reefs. The concrete creates shelter & a home for marine life.
See also
Eurocode 2: Design of concrete structures
References
Concrete | Types of concrete | [
"Engineering"
] | 5,990 | [
"Structural engineering",
"Concrete"
] |
24,969,353 | https://en.wikipedia.org/wiki/Plasma%20lamp | Plasma lamps are a type of electrodeless gas-discharge lamp energized by radio frequency (RF) power. They are distinct from the novelty plasma lamps that were popular in the 1980s.
The internal-electrodeless lamp was invented by Nikola Tesla after his experimentation with high-frequency currents in evacuated glass tubes for the purposes of lighting and the study of high voltage phenomena. The first practical plasma lamps were the sulfur lamps manufactured by Fusion Lighting. This lamp suffered several practical problems and did not prosper commercially. Plasma lamps with an internal phosphor coating are called external electrode fluorescent lamps (EEFL); these external electrodes or terminal conductors provide the radio frequency electric field.
Description
Modern plasma lamps are a family of light sources that generate light by exciting plasma inside a closed transparent burner or bulb using radio frequency (RF) power. Typically, such lamps use a noble gas or a mixture of these gases and additional materials such as metal halides, sodium, mercury or sulfur. In modern plasma lamps, a waveguide is used to constrain and focus the electrical field into the plasma. In operation, the gas is ionized, and free electrons, accelerated by the electrical field, collide with gas and metal atoms. Some atomic electrons circling around the gas and metal atoms are excited by these collisions, bringing them to a higher energy state. When the electron falls back to its original state, it emits a photon, resulting in visible light or ultraviolet radiation, depending on the fill materials.
The first commercial plasma lamp was an ultraviolet curing lamp with a bulb filled with argon and mercury vapor developed by Fusion UV. That lamp led Fusion Lighting to the development of the sulfur lamp, a bulb filled with argon and sulfur that is bombarded with microwaves through a hollow waveguide. The bulb had to be spun rapidly to prevent the sulfur from burning through. Fusion Lighting did not prosper commercially, but other manufacturers continue to pursue sulfur lamps. Sulfur lamps, though relatively efficient, have had several problems, chiefly:
Limited life – Magnetrons had limited lives.
Large size
Heat – The sulfur burnt through the bulb wall unless it was rotated rapidly.
High power demand – They were not able to sustain a plasma in powers under 1000 W.
Limited life
In the past, the life of the plasma lamps was limited by the magnetron used to generate the microwaves. Solid-state RF chips can be used and give long lives. However, using solid-state chips to generate RF is currently an order of magnitude more expensive than using a magnetron and so only appropriate for high-value lighting niches. It has recently been shown by Dipolar of Sweden to be possible to extend the life of magnetrons to over 40,000 hours, making low-cost plasma lamps possible.
Heat and power
The use of a high-dielectric waveguide allowed the sustaining of plasmas at much lower powers—down to 100 W in some instances. It also allowed the use of conventional gas-discharge lamp fill materials which removed the need to spin the bulb. The only issue with the ceramic waveguide was that much of the light generated by the plasma was trapped inside the opaque ceramic waveguide.
High-efficiency plasma (HEP)
High-efficiency plasma lighting is the class of plasma lamps that have system efficiencies of 90 lumens per watt or more. Lamps in this class are potentially the most energy-efficient light source for outdoor, commercial, and industrial lighting. This is due not only to their high system efficiency but also to the small light source they present enabling very high luminaire efficiency.
Luminaire Efficacy Rating (LER) is the single figure of merit the National Electrical Manufacturers Association has defined to help address problems with lighting manufacturers' efficiency claims and is designed to allow robust comparison between lighting types. It is given by the product of luminaire efficiency (EFF) times total rated lamp output in lumens (TLL) times ballast factor (BF), divided by the input power in watts (IP):
LER = EFF × TLL × BF / IP
The "system efficiency" for a high-efficiency plasma lamp is given by the last three variables, that is, it excludes the luminaire efficiency. Though plasma lamps do not have a ballast, they have an RF power supply that fulfills the equivalent function. In electrodeless lamps, the inclusion of the electrical losses, or "ballast factor", in lumens per watt claimed can be particularly significant as the conversion of electrical power to radio frequency (RF) power can be a highly inefficient process.
Many modern plasma lamps have very small light sources—far smaller than HID bulbs or fluorescent tubes—leading to much higher luminaire efficiencies also. High-intensity discharge lamps have typical luminaire efficiencies of 55%, and fluorescent lamps of 70%. Plasma lamps typically have luminaire efficiencies exceeding 90%.
Applications
Plasma lamps have been used in high bay and street lighting applications, as well as in stage lighting. They were briefly used in some projection televisions.
References
Gas discharge lamps
Types of lamp
Plasma technology and applications | Plasma lamp | [
"Physics"
] | 1,051 | [
"Plasma technology and applications",
"Plasma physics"
] |
23,467,710 | https://en.wikipedia.org/wiki/Exponential%20map%20%28discrete%20dynamical%20systems%29 | In the theory of dynamical systems, the exponential map can be used as the evolution function of the discrete nonlinear dynamical system.
Family
The family of exponential functions is called the exponential family.
Forms
There are many forms of these maps, many of which are equivalent under a coordinate transformation. For example two of the most common ones are:
The second one can be mapped to the first using the fact that , so is the same under the transformation . The only difference is that, due to multi-valued properties of exponentiation, there may be a few select cases that can only be found in one version. Similar arguments can be made for many other formulas.
References
Chaotic maps | Exponential map (discrete dynamical systems) | [
"Mathematics"
] | 137 | [
"Mathematical analysis",
"Functions and mappings",
"Mathematical analysis stubs",
"Mathematical objects",
"Mathematical relations",
"Chaotic maps",
"Dynamical systems"
] |
23,470,492 | https://en.wikipedia.org/wiki/Bidomain%20model | The bidomain model is a mathematical model to define the electrical activity of the heart. It consists in a continuum (volume-average) approach in which the cardiac microstructure is defined in terms of muscle fibers grouped in sheets, creating a complex three-dimensional structure with anisotropical properties. Then, to define the electrical activity, two interpenetrating domains are considered, which are the intracellular and extracellular domains, representing respectively the space inside the cells and the region between them.
The bidomain model was first proposed by Schmitt in 1969 before being formulated mathematically in the late 1970s.
Since it is a continuum model, rather than describing each cell individually, it represents the average properties and behaviour of group of cells organized in complex structure. Thus, the model results to be a complex one and can be seen as a generalization of the cable theory to higher dimensions and, going to define the so-called bidomain equations.
Many of the interesting properties of the bidomain model arise from the condition of unequal anisotropy ratios. The electrical conductivity in anisotropic tissues is not unique in all directions, but it is different in parallel and perpendicular direction with respect to the fiber one.
Moreover, in tissues with unequal anisotropy ratios, the ratio of conductivities parallel and perpendicular to the fibers are different in the intracellular and extracellular spaces. For instance, in cardiac tissue, the anisotropy ratio in the intracellular space is about 10:1, while in the extracellular space it is about 5:2.
Mathematically, unequal anisotropy ratios means that the effect of anisotropy cannot be removed by a change in the distance scale in one direction.
Instead, the anisotropy has a more profound influence on the electrical behavior.
Three examples of the impact of unequal anisotropy ratios are
the distribution of transmembrane potential during unipolar stimulation of a sheet of cardiac tissue,
the magnetic field produced by an action potential wave front propagating through cardiac tissue,
the effect of fiber curvature on the transmembrane potential distribution during an electric shock.
Formulation
Bidomain domain
The bidomain domain is principally represented by two main regions: the cardiac cells, called intracellular domain, and the space surrounding them, called extracellular domain. Moreover, usually another region is considered, called extramyocardial region.
The intracellular and extracellular domains, which are separate by the cellular membrane, are considered to be a unique physical space representing the heart (), while the extramyocardial domain is a unique physical space adjacent of them (). The extramyocardial region can be considered as a fluid bath, especially when one wants to simulate experimental conditions, or as a human torso to simulate physiological conditions.
The boundary of the two principal physical domains defined are important to solve the bidomain model. Here the heart boundary is denoted as while the torso domain boundary is
Unknowns and parameters
The unknowns in the bidomain model are three, the intracellular potential , the extracellular potential and the transmembrane potential , which is defined as the difference of the potential across the cell membrane .
Moreover, some important parameters need to be taken in account, especially the intracellular conductivity tensor matrix the extracellular conductivity tensor matrix The transmembrane current flows between the intracellular and extracellular regions and it is in part described by the corresponding ionic current over the membrane per unit area .
Moreover, the membrane capacitance per unit area and the surface to volume ratio of the cell membrane need to be considered to derive the bidomain model formulation, which is done in the following section.
Standard formulation
The bidomain model is defined through two partial differential equations (PDE) the first of which is a reaction diffusion equation in terms of the transmembrane potential, while the second one computes the extracellular potential starting from a given transmembran potential distribution.
Thus, the bidomain model can be formulated as follows:
where and can be defined as applied external stimulus currents.
Ionic current equation
The ionic current is usually represented by an ionic model through a system of ordinary differential equations (ODEs). Mathematically, one can write where is called ionic variable. Then, in general, for all , the system reads
Different ionic models have been proposed:
phenomenological models, which are the simplest ones and used to reproduce macroscopic behavior of the cell.
physiological models, which take into account both macroscopic behaviour and cell physiology with a quite detailed description of the most important ionic current.
Model of an extramyocardial region
In some cases, an extramyocardial region is considered. This implies the addition to the bidomain model of an equation describing the potential propagation inside the extramyocardial domain.
Usually, this equation is a simple generalized Laplace equation of type
where is the potential in the extramyocardial region and is the corresponding conductivity tensor.
Moreover, an isolated domain assumption is considered, which means that the following boundary conditions are added
being the unit normal directed outside of the extramyocardial domain.
If the extramyocardial region is the human torso, this model gives rise to the forward problem of electrocardiology.
Derivation
The bidomain equations are derived from the Maxwell's equations of the electromagnetism, considering some simplifications.
The first assumption is that the intracellular current can flow only between the intracellular and extracellular regions, while the intracellular and extramyocardial regions can comunicate between them, so that the current can flow into and from the extramyocardial regions but only in the extracellular space.
Using Ohm's law and a quasi-static assumption, the gradient of a scalar potential field can describe an electrical field , which means that
Then, if represent the current density of the electric field , two equations can be obtained
where the subscript and represent the intracellular and extracellular quantities respectively.
The second assumption is that the heart is isolated so that the current that leaves one region need to flow into the other. Then, the current density in each of the intracellular and extracellular domain must be equal in magnitude but opposite in sign, and can be defined as the product of the surface to volume ratio of the cell membrane and the transmembrane ionic current density per unit area, which means that
By combining the previous assumptions, the conservation of current densities is obtained, namely
from which, summing the two equations
This equation states exactly that all currents exiting one domain must enter the other.
From here, it is easy to find the second equation of the bidomain model subtracting from both sides. In fact,
and knowing that the transmembral potential is defined as
Then, knowing the transmembral potential, one can recover the extracellular potential.
Then, the current that flows across the cell membrane can be modelled with the cable equation,
Combining equations () and () gives
Finally, adding and subtracting on the left and rearranging , one can get the first equation of the bidomain model
which describes the evolution of the transmembrane potential in time.
The final formulation described in the standard formulation section is obtained through a generalization, considering possible external stimulus which can be given through the external applied currents and .
Boundary conditions
In order to solve the model, boundary conditions are needed. The more classical boundary conditions are the following ones, formulated by Tung.
First of all, as state before in the derive section, there ca not been any flow of current between the intracellular and extramyocardial domains. This can be mathematically described as
where is the vector that represents the outwardly unit normal to the myocardial surface of the heart.
Since the intracellular potential is not explicitily presented in the bidomain formulation, this condition is usually described in terms of the transmembrane and extracellular potential, knowing that , namely
For the extracellular potential, if the myocardial region is presented, a balance in the flow between the extracellular and the extramyocardial regions is considered
Here the normal vectors from the perspective of both domains are considered, thus the negative sign are necessary. Moreover, a perfect transmission of the potential on the cardiac boundary is necessary, which gives
Instead, if the heart is considered as isolated, which means that no myocardial region is presented, a possible boundary condition for the extracellular problem is
Reduction to monodomain model
By assuming equal anisotropy ratios for the intra- and extracellular domains, i.e. for some scalar , the model can be reduced to one single equation, called monodomain equation
where the only variable is now the transmembrane potential, and the conductivity tensor is a combination of and
Formulation with boundary conditions in an isolated domain
If the heart is considered as an isolated tissue, which means that no current can flow outside of it, the final formulation with boundary conditions reads
Numerical solution
There are various possible techniques to solve the bidomain equations. Between them, one can find finite difference schemes, finite element schemes and also finite volume schemes. Special considerations can be made for the numerical solution of these equations, due to the high time and space resolution needed for numerical convergence.
See also
Monodomain model
Forward problem of electrocardiology
References
External links
Scholarpedia article about the bidomain model
Cardiac electrophysiology
Electrophysiology
Partial differential equations
Mathematical modeling
Numerical analysis | Bidomain model | [
"Mathematics"
] | 1,941 | [
"Mathematical modeling",
"Applied mathematics",
"Computational mathematics",
"Mathematical relations",
"Numerical analysis",
"Approximations"
] |
23,470,812 | https://en.wikipedia.org/wiki/Biotechnology%20and%20Bioprocess%20Engineering | Biotechnology and Bioprocess Engineering is a peer-reviewed bimonthly scientific journal published by Springer Science+Business Media on behalf of the Korean Society for Biotechnology and Bioengineering. Biotechnology and Bioprocess Engineering covers all aspects of biotechnology and bioengineering. The editor-in-chief of the journal is Jong Won Yun (Daegu University) and Sang Yup Lee (KAIST). The founding editors-in-chief were Cha-Yong Choi (Seoul National University), Ho Nam Chang (Korea Advanced Institute of Science and Technology), and Sun Bok Lee (POSTECH).
Abstracting and indexing
The journal is abstracted and indexed in EMBASE, EMBiology, Food Science and Technology Abstracts, IBIDS, International Abstracts in Operations Research, Science Citation Index Expanded, and Scopus. According to the Journal Citation Reports, its 2021 impact factor is 3.386
References
External links
Biotechnology journals
Bimonthly journals
Springer Science+Business Media academic journals
English-language journals
Academic journals established in 1996 | Biotechnology and Bioprocess Engineering | [
"Biology"
] | 212 | [
"Biotechnology literature",
"Biotechnology journals"
] |
23,471,261 | https://en.wikipedia.org/wiki/Absorptance | In the study of heat transfer, absorptance of the surface of a material is its effectiveness in absorbing radiant energy. It is the ratio of the absorbed to the incident radiant power.
Mathematical definitions
Hemispherical absorptance
Hemispherical absorptance of a surface, denoted is defined as
where
is the radiant flux absorbed by that surface;
is the radiant flux received by that surface.
Spectral hemispherical absorptance
Spectral hemispherical absorptance in frequency and spectral hemispherical absorptance in wavelength of a surface, denoted and respectively, are defined as
where
is the spectral radiant flux in frequency absorbed by that surface;
is the spectral radiant flux in frequency received by that surface;
is the spectral radiant flux in wavelength absorbed by that surface;
is the spectral radiant flux in wavelength received by that surface.
Directional absorptance
Directional absorptance of a surface, denoted , is defined as
where
is the radiance absorbed by that surface;
is the radiance received by that surface.
Spectral directional absorptance
Spectral directional absorptance in frequency and spectral directional absorptance in wavelength of a surface, denoted and respectively, are defined as
where
is the spectral radiance in frequency absorbed by that surface;
is the spectral radiance received by that surface;
is the spectral radiance in wavelength absorbed by that surface;
is the spectral radiance in wavelength received by that surface.
Other radiometric coefficients
References
Physical quantities
Radiometry | Absorptance | [
"Physics",
"Mathematics",
"Engineering"
] | 311 | [
"Physical phenomena",
"Telecommunications engineering",
"Physical quantities",
"Quantity",
"Physical properties",
"Radiometry"
] |
23,473,595 | https://en.wikipedia.org/wiki/Light-year | A light-year, alternatively spelled light year (ly or lyr), is a unit of length used to express astronomical distances and is equal to exactly , which is approximately 5.88 trillion mi. As defined by the International Astronomical Union (IAU), a light-year is the distance that light travels in vacuum in one Julian year (365.25 days). Despite its inclusion of the word "year", the term should not be misinterpreted as a unit of time.
The light-year is most often used when expressing distances to stars and other distances on a galactic scale, especially in non-specialist contexts and popular science publications. The unit most commonly used in professional astronomy is the parsec (symbol: pc, about 3.26 light-years).
Definitions
As defined by the International Astronomical Union (IAU), the light-year is the product of the Julian year (365.25 days, as opposed to the 365.2425-day Gregorian year or the 365.24219-day Tropical year that both approximate) and the speed of light (). Both of these values are included in the IAU (1976) System of Astronomical Constants, used since 1984. From this, the following conversions can be derived:
{|
|-
|rowspan=6 valign=top|1 light-year
|= metres (exactly)
|-
|≈ petametres
|-
|≈ trillion (short scale) kilometres ( trillion miles)
|-
|≈ astronomical units
|-
|≈ parsec
|}
The abbreviation used by the IAU for light-year is "ly", International standards like ISO 80000:2006 (now superseded) have used "l.y." and localized abbreviations are frequent, such as "al" in French, Spanish, and Italian (from année-lumière, año luz and anno luce, respectively), "Lj" in German (from Lichtjahr), etc.
Before 1984, the tropical year (not the Julian year) and a measured (not defined) speed of light were included in the IAU (1964) System of Astronomical Constants, used from 1968 to 1983. The product of Simon Newcomb's J1900.0 mean tropical year of ephemeris seconds and a speed of light of produced a light-year of (rounded to the seven significant digits in the speed of light) found in several modern sources was probably derived from an old source such as C. W. Allen's 1973 Astrophysical Quantities reference work, which was updated in 2000, including the IAU (1976) value cited above (truncated to 10 significant digits).
Other high-precision values are not derived from a coherent IAU system. A value of found in some modern sources is the product of a mean Gregorian year (365.2425 days or ) and the defined speed of light (). Another value, , is the product of the J1900.0 mean tropical year and the defined speed of light.
Abbreviations used for light-years and multiples of light-years are:
"ly" for one light-year
"kly" or "klyr" for a kilolight-year (1,000 light-years)
"Mly" for a megalight-year (1,000,000 light-years)
"Gly" or "Glyr" for a gigalight-year ( light-years)
History
The light-year unit appeared a few years after the first successful measurement of the distance to a star other than the Sun, by Friedrich Bessel in 1838. The star was 61 Cygni, and he used a heliometre designed by Joseph von Fraunhofer. The largest unit for expressing distances across space at that time was the astronomical unit, equal to the radius of the Earth's orbit at . In those terms, trigonometric calculations based on 61 Cygni's parallax of 0.314 arcseconds, showed the distance to the star to be . Bessel added that light takes 10.3 years to traverse this distance. He recognized that his readers would enjoy the mental picture of the approximate transit time for light, but he refrained from using the light-year as a unit. He may have resisted expressing distances in light-years because it would reduce the accuracy of his parallax data due to multiplying with the uncertain parameter of the speed of light.
The speed of light was not yet precisely known in 1838; the estimate of its value changed in 1849 (Fizeau) and 1862 (Foucault). It was not yet considered to be a fundamental constant of nature, and the propagation of light through the aether or space was still enigmatic.
The light-year unit appeared in 1851 in a German popular astronomical article by Otto Ule. Ule explained the oddity of a distance unit name ending in "year" by comparing it to a walking hour (Wegstunde).
A contemporary German popular astronomical book also noticed that light-year is an odd name. In 1868 an English journal labelled the light-year as a unit used by the Germans. Eddington called the light-year an inconvenient and irrelevant unit, which had sometimes crept from popular use into technical investigations.
Although modern astronomers often prefer to use the parsec, light-years are also popularly used to gauge the expanses of interstellar and intergalactic space.
Usage of term
Distances expressed in light-years include those between stars in the same general area, such as those belonging to the same spiral arm or globular cluster. Galaxies themselves span from a few thousand to a few hundred thousand light-years in diameter, and are separated from neighbouring galaxies and galaxy clusters by millions of light-years. Distances to objects such as quasars and the Sloan Great Wall run up into the billions of light-years.
Related units
Distances between objects within a star system tend to be small fractions of a light-year, and are usually expressed in astronomical units. However, smaller units of length can similarly be formed usefully by multiplying units of time by the speed of light. For example, the light-second, useful in astronomy, telecommunications and relativistic physics, is exactly metres or of a light-year. Units such as the light-minute, light-hour and light-day are sometimes used in popular science publications. The light-month, roughly one-twelfth of a light-year, is also used occasionally for approximate measures. The Hayden Planetarium specifies the light month more precisely as 30 days of light travel time.
Light travels approximately one foot in a nanosecond; the term "light-foot" is sometimes used as an informal measure of time.
See also
1 petametre (examples of distances on the order of one light-year)
Einstein protocol
Hubble length
Orders of magnitude (length)
Notes
References
External links
Light
Units of length
Units of measurement in astronomy
Concepts in astronomy
1838 in science | Light-year | [
"Physics",
"Astronomy",
"Mathematics"
] | 1,448 | [
"Physical phenomena",
"Spectrum (physical sciences)",
"Units of length",
"Concepts in astronomy",
"Quantity",
"Electromagnetic spectrum",
"Units of measurement in astronomy",
"Waves",
"Light",
"Units of measurement"
] |
23,473,596 | https://en.wikipedia.org/wiki/Electron%20wake | Electron wake is the disturbance left after a high-energy charged particle passes through condensed matter or plasma. Ions passing through can introduce periodic oscillations in the crystal lattice or plasma wave with the characteristic frequency of the crystal or plasma frequency. Interactions of the field created by these oscillations with the charged particle field alternate from constructive interference to destructive interference, producing alternating waves of electric field and displacement. The frequency of the wake field is determined by the nature of the penetrated matter, and the period of the wake field is directly proportional to the speed of the incoming charged particle. The amplitude of the first wake wave is the most important, as it produces a braking force on the charged particle, eventually slowing it down. Wake fields also can capture and guide lightweight ions or positrons in the direction perpendicular to the wake. The larger the speed of the original charged particle, the larger the angle between the initial particle's velocity and the captured ion's velocity.
References
Encyclopedia article on the electron wake
On the possibility of accelerating positron on an electron wake at SABER
See also
Coulomb explosion
Charged particle beam
Wake fields
Plasma acceleration
Bremsstrahlung
Atomic physics
Plasma phenomena
Scattering
Accelerator physics | Electron wake | [
"Physics",
"Chemistry",
"Materials_science"
] | 243 | [
"Physical phenomena",
"Applied and interdisciplinary physics",
"Plasma physics",
"Plasma phenomena",
"Scattering stubs",
"Quantum mechanics",
"Scattering",
"Plasma physics stubs",
"Particle physics",
" molecular",
"Nuclear physics",
"Atomic",
"Atomic physics",
"Condensed matter physics",
... |
23,475,491 | https://en.wikipedia.org/wiki/Luttinger%E2%80%93Kohn%20model | The Luttinger–Kohn model is a flavor of the k·p perturbation theory used for calculating the structure of multiple, degenerate electronic bands in bulk and quantum well semiconductors. The method is a generalization of the single band k·p theory.
In this model, the influence of all other bands is taken into account by using Löwdin's perturbation method.
Background
All bands can be subdivided into two classes:
Class A: six valence bands (heavy hole, light hole, split off band and their spin counterparts) and two conduction bands.
Class B: all other bands.
The method concentrates on the bands in Class A, and takes into account Class B bands perturbatively.
We can write the perturbed solution, , as a linear combination of the unperturbed eigenstates :
Assuming the unperturbed eigenstates are orthonormalized, the eigenequations are:
,
where
.
From this expression, we can write:
,
where the first sum on the right-hand side is over the states in class A only, while the second sum is over the states on class B. Since we are interested in the coefficients for m in class A, we may eliminate those in class B by an iteration procedure to obtain:
,
Equivalently, for ():
and
.
When the coefficients belonging to Class A are determined, so are .
Schrödinger equation and basis functions
The Hamiltonian including the spin-orbit interaction can be written as:
,
where is the Pauli spin matrix vector. Substituting into the Schrödinger equation in Bloch approximation we obtain
,
where
and the perturbation Hamiltonian can be defined as
The unperturbed Hamiltonian refers to the band-edge spin-orbit system (for k=0). At the band edge, the conduction band Bloch waves exhibits s-like symmetry, while the valence band states are p-like (3-fold degenerate without spin). Let us denote these states as , and , and respectively. These Bloch functions can be pictured as periodic repetition of atomic orbitals, repeated at intervals corresponding to the lattice spacing. The Bloch function can be expanded in the following manner:
,
where j' is in Class A and is in Class B. The basis functions can be chosen to be
.
Using Löwdin's method, only the following eigenvalue problem needs to be solved
where
,
The second term of can be neglected compared to the similar term with p instead of k. Similarly to the single band case, we can write for
We now define the following parameters
and the band structure parameters (or the Luttinger parameters) can be defined to be
These parameters are very closely related to the effective masses of the holes in various valence bands. and describe the coupling of the , and states to the other states. The third parameter relates to the anisotropy of the energy band structure around the point when .
Explicit Hamiltonian matrix
The Luttinger-Kohn Hamiltonian can be written explicitly as a 8X8 matrix (taking into account 8 bands - 2 conduction, 2 heavy-holes, 2 light-holes and 2 split-off)
Summary
References
2. Luttinger, J. M. Kohn, W., "Motion of Electrons and Holes in Perturbed Periodic Fields", Phys. Rev. 97,4. pp. 869-883, (1955). https://journals.aps.org/pr/abstract/10.1103/PhysRev.97.869
Condensed matter physics | Luttinger–Kohn model | [
"Physics",
"Chemistry",
"Materials_science",
"Engineering"
] | 750 | [
"Phases of matter",
"Condensed matter physics",
"Matter",
"Materials science"
] |
23,475,732 | https://en.wikipedia.org/wiki/Local%20oxidation%20nanolithography | Local oxidation nanolithography (LON) is a tip-based nanofabrication method. It is based on the spatial confinement on an oxidation reaction under the sharp tip of an atomic force microscope.
The first materials on which LON was demonstrated were Si(111) and polycrystalline tantalum. Subsequently, the technique has been extended to III–V semiconductors, silicon carbide, metals such as titanium, tantalum, aluminium, molybdenum, nickel and niobium; thin films of manganite in the perovskite form; dielectrics like silicon nitride, organosilane self-assembled monolayers, dendritic macromolecules and carbonaceous films.
History
The local oxidation of a surface by means of a scanning probe technique was first observed by Dagata and co-workers in 1990 who locally modified a hydrogen-terminated silicon surface into silicon dioxide by applying a bias voltage between the tip of a scanning tunneling microscope and the surface itself. In 1993 Day and Allee demonstrated the possibility of performing local oxidation experiments with an atomic force microscope, which opened the way to applying the technique to a large variety of materials.
Basic principle
Currently, local oxidation experiments are performed with an atomic force microscope operated in contact or noncontact mode with additional circuits to apply voltage pulses between tip and sample. The local oxidation process is mediated by the formation of a water meniscus.
In order to perform local oxidation nanolithography, the relative humidity in the AFM chamber is kept between 30% and 60%. A voltage pulse is applied between a conductive AFM tip and the sample. The applied voltage induces the formation of a water bridge between tip and sample whenever the amplitude of the voltage pulse is above a certain threshold voltage. When the liquid meniscus is created the applied voltage pulse causes an oxidation reaction by breaking the covalent bonds in the water molecules. The liquid bridge provides the oxyanions (OH−,O−) needed to form the oxide and confines the lateral extension of the region to be oxidized.
The chemical reactions that govern the local oxidation in a metallic substrate (M) are the following:
while hydrogen gas is liberated at the AFM tip through the reduction reaction:
2H+ + 2e- -> H2
When the voltage pulse is off the AFM feedback forces the cantilever to recover its original oscillation amplitude withdrawing the tip from the sample and breaking the liquid meniscus. Finally the AFM continues to scan the sample thus allowing to image MOn nanostructure fabricated during the Local Oxidation process with the very same tip used for its fabrication.
The method to form liquid bridges is so precise that water meniscus diameters of 20 nm or below are easily obtained. This has led to the reproducible fabrication of sub-10 nm structures in silicon and other metallic surfaces.
Experimental setup
Local oxidation experiments can be performed with almost any kind of atomic force microscope. The key requirement is the possibility to apply voltage pulses between the tip and the sample. It is recommendable to enclose the microscope in a chamber where the atmosphere is controlled. In the simplest case, the oxidant is water vapor, which is naturally present in the air. Controlling the relative humidity generally helps to obtain more reproducible results. The size of the fabricated features depends on a number of parameters, such as the distance between the sample and the tip, the amplitude and the duration of the voltage pulse, and the relative humidity of the atmosphere.
Applications
The development of nanometer-scale lithographies is the focus of an intense research activity because progress on nanotechnology depends on the capability to fabricate, position and interconnect nanometer-scale structures.
Patterning
Local Oxidation Nanolithography allows to create a large variety of motives like dots, lines and letters with nanometer accuracy. In 2005, researchers at the Spanish National Research Council in Madrid wrote the first ten lines of Cervantes' Don Quixote on a few square micrometres of silicon. This pattern versatility can be used for information storage or to design etch-resistant nanomasks in order to fabricate nanodevices as well as many other applications.
Data storage
It is possible to store information using dot-like nanostructures created by the local oxidation of a surface. This storage uses the binary code considering the presence of a nanostructure as a 1 and its lack as a 0. In this way information can be stored in a small surface with a single SiO2 dot constituting a bit. In 1999 Cooper et al. demonstrated that this methods allows to obtain an information density of 1.6 Tbit/in2. However, only read-only memories can be fabricated with this technique.
Molecular template growth and preferential deposition
Local oxidation of silicon surfaces by noncontact atomic-force microscopy is an emerging and promising method for patterning surfaces at the nanometer scale due to its very precise control of the feature size. The features created with this technique can be used for the template growth and preferential deposition of different molecules like single-molecule magnets, biomolecules and conjugated organic molecules.
This method of nanopositioning is an important tool for the fabrication of new nanodevices based on the novel properties exhibited by some nanoparticles and molecules. Potential applications of single-molecule magnets (SMMs) such as Mn12 as bits for information storage or qubits for quantum computation require methods for nanoscale-controlled positioning and/or manipulation of those molecules. The patterning of the Mn12 molecules on a silicon surface is achieved by first derivatizing this surface with a self-assembled monolayer of APTES, which leaves it terminated by amino groups (-NH2). Such termination electrostatically repels the Mn12 molecules. Subsequently, a pattern of silicon dioxide is defined by LON. The SMM molecules are predominantly deposited on the oxide motives because of electrostatic attraction. The electrostatic attraction between the silicon oxide fabricated by LON and the Mn12 molecules achieves the preferential deposition of this molecules with a nanoscale accuracy.
Fabrication of nanodevices
By using local oxidation nanolithography as tool for the fabrication of etch-resistant nanomasks, it is possible to fabricate nanoscale electronic devices, such as field-effect transistors, single-electron transistors, Josephson junctions, quantum rings or SQUIDs. LON also allows to fabricate silicon nanowires (SiNWs) in a top-down fashion starting from silicon on insulator (SOI) wafers. Local oxidation nanolithography contributes to the nanometric precision of the device fabrication. This top-down fabrication technique allows the fabrication of a large variety of SiNWs with different shapes, from angular to circular. It also allows the precise positioning of the silicon nanowires in any desired position, making easier its integration; indeed, this technique is compatible with the standard silicon CMOS processing technology. Single crystalline silicon nanowires have already shown a great potential as ultrasensitive sensors by detecting changes in the nanowire conductivity when a specific analyte is present. Local oxidation nanolithography, therefore, is a promising technique to allow the realisation of array of biosensors.
References
External links
Local oxidation nanolithography page in García's research group at CSIC
Miles' research group at the University of Bristol
Quate's group at Stanford University
Scanning probe microscopy
Lithography (microfabrication) | Local oxidation nanolithography | [
"Chemistry",
"Materials_science"
] | 1,557 | [
"Microtechnology",
"Scanning probe microscopy",
"Microscopy",
"Nanotechnology",
"Lithography (microfabrication)"
] |
23,476,130 | https://en.wikipedia.org/wiki/Neutrino%20decoupling | In Big Bang cosmology, neutrino decoupling was the epoch at which neutrinos ceased interacting with other types of matter, and thereby ceased influencing the dynamics of the universe at early times. Prior to decoupling, neutrinos were in thermal equilibrium with protons, neutrons and electrons, which was maintained through the weak interaction. Decoupling occurred approximately at the time when the rate of those weak interactions was slower than the rate of expansion of the universe. Alternatively, it was the time when the time scale for weak interactions became greater than the age of the universe at that time. Neutrino decoupling took place approximately one second after the Big Bang, when the temperature of the universe was approximately 10 billion kelvin, or 1 MeV.
As neutrinos rarely interact with matter, these neutrinos still exist today, analogous to the much later cosmic microwave background emitted during recombination, around 377,000 years after the Big Bang. They form the cosmic neutrino background (abbreviated CνB or CNB). The neutrinos from this event have a very low energy, around 10−10 times smaller than is possible with present-day direct detection. Even high energy neutrinos are notoriously difficult to detect, so the CNB may not be directly observed in detail for many years, if at all. However, Big Bang cosmology makes many predictions about the CNB, and there is very strong indirect evidence that the CNB exists.
Derivation of decoupling time
Neutrinos are scattered (interfering with free streaming) by their interactions with electrons and positrons, such as the reaction
.
The approximate rate of these interactions is set by the number density of electrons and positrons, the averaged product of the cross section for interaction and the velocity of the particles. The number density of the relativistic electrons and positrons depends on the cube of the temperature , so that . The product of the cross section and velocity for weak interactions for temperatures (energies) below W/Z boson masses (~100 GeV) is given approximately by , where is Fermi's constant (as is standard in particle physics calculations, factors of the speed of light are set equal to 1). Putting it all together, the rate of weak interactions is
.
This can be compared to the expansion rate which is given by the Hubble parameter , with
,
where is the gravitational constant and is the energy density of the universe. At this point in cosmic history, the energy density is dominated by radiation, so that . As the rate of weak interaction depends more strongly on temperature, it will fall more quickly as the universe cools. Thus when the two rates are approximately equal (dropping terms of order unity, including an effective degeneracy term which counts the number of states of particles which are interacting) gives the approximate temperature at which neutrinos decouple:
.
Solving for temperature gives
.
While this is a very rough derivation, it illustrates the important physical phenomena which determined when neutrinos decoupled.
Observational evidence
While neutrino decoupling can not be observed directly, it is expected to have left behind a cosmic neutrino background, analogous to the cosmic microwave background radiation of electromagnetic radiation which was emitted at a much later epoch. "The detection of the neutrino background is far beyond the capabilities of the present generation of neutrino detectors." There is data, however, which indirectly indicates the presence of a neutrino background. One piece of evidence is damping of the angular power spectrum of the CMB, which results from anisotropies in the neutrino background.
Another indirect measurement of neutrino decoupling is allowed by the role that neutrino decoupling plays in setting the ratio of neutrons to protons. Before decoupling, the number of neutrons and protons are maintained in their equilibrium abundances by weak interactions, specifically beta decay and electron capture (or inverse beta decay) according to
and
.
Once the rate of weak interactions is slower than the characteristic rate of the expansion of the universe, this equilibrium cannot be maintained, and the abundance of neutrons to protons "freezes in," at a value
.
This value is simply found by evaluating the Boltzmann factor for neutrons and protons at decoupling time, according to
,
where is the mass difference between neutrons and protons and is the temperature at decoupling. This ratio is critical to the synthesis of atoms during Big Bang nucleosynthesis, the process which formed the majority of helium atoms in the universe, as it "is the dominant factor in determining the amount of helium produced." As helium atoms are stable, the neutrons are locked in, and beta decay of neutrons into protons, electrons, and neutrinos can no longer occur. Thus the abundance of neutrons in the primordial matter can be measured by astronomers, and, as it was determined by the ratio of neutrons to protons at neutrino decoupling, the helium abundance indirectly measures the temperature at which neutrino decoupling took place, and is in agreement with the figure derived above.
Indirect evidence from phase changes to the Cosmic Microwave Background (CMB)
Big Bang cosmology makes many predictions about the CNB, and there is very strong indirect evidence that the cosmic neutrino background exists, both from Big Bang nucleosynthesis predictions of the helium abundance, and from anisotropies in the cosmic microwave background. One of these predictions is that neutrinos will have left a subtle imprint on the cosmic microwave background (CMB). It is well known that the CMB has irregularities. Some of the CMB fluctuations were roughly regularly spaced, because of the effect of baryon acoustic oscillations. In theory, the decoupled neutrinos should have had a very slight effect on the phase of the various CMB fluctuations.
In 2015, it was reported that such shifts had been detected in the CMB. Moreover, the fluctuations corresponded to neutrinos of almost exactly the temperature predicted by Big Bang theory ( compared to a prediction of 1.95 K), and exactly three types of neutrino, the same number of neutrino flavours currently predicted by the Standard Model.
See also
Chronology of the universe
References
Bibliography
External links
Student notes from a UC Berkeley cosmology course (see page 16)
Professor notes from a University of Oregon astrophysics course
Physical cosmology | Neutrino decoupling | [
"Physics",
"Astronomy"
] | 1,347 | [
"Astronomical sub-disciplines",
"Theoretical physics",
"Physical cosmology",
"Astrophysics"
] |
23,476,165 | https://en.wikipedia.org/wiki/Sweep%20generator | A sweep generator is a piece of electronic test equipment similar to, and sometimes included on, a function generator which creates an electrical waveform with a linearly varying frequency and a constant amplitude. Sweep generators are commonly used to test the frequency response of electronic filter circuits. These circuits are mostly transistor circuits with inductors and capacitors to create linear characteristics.
Sweeps are a popular method in the field of audio measurement to describe the change in a measured output value over a progressing input parameter. The most commonly-used progressive input parameter is frequency varied over the standard audio bandwidth of 20 Hz to 20 kHz.
Glide Sweep
A glide sweep (or chirp) is a continuous signal in which the frequency increases or decreases logarithmically with time. This provides the complete range of testing frequencies between the start and stop frequency. An advantage over the stepped sweep is that the signal duration can be reduced by the user without any loss of frequency resolution in the results. This allows for rapid testing. Although the theory behind the glide sweep has been known for several decades, its use in audio measuring devices has only evolved over the past several years. The reason for this lies with the high computing power required.
Stepped Sweep
In a stepped sweep, one variable input parameter (frequency or amplitude) is incremented or decremented in discrete steps. After each change, the analyzer waits until a stable reading is detected before switching to the next step. The scaling of the steps is linear or logarithmic. Since the settling time of different test objects cannot be predicted, the duration of a stepped sweep cannot be determined exactly in advance. For the determination of amplitude or frequency response, the stepped sweep has been largely replaced by the glide sweep. The main application for the stepped sweep is to measure the linearity of systems. Here, the frequency of the test signal is kept constant while the amplitude is varied. Typically the amplitude and distortion of the device under test are measured. This is also referred to as an "amplitude sweep".
Time Sweep
In the case of a time sweep, the x-axis represents time. Again the y-axis represents a measured value, e.g. amplitude. The change in the measured value is observed over time. For example, how does the response of the device under test change over a long period?
Table Sweep
A rarely used special form of the stepped sweep is the table sweep. Here the input signal is produced from a table as a sequence of any frequency and amplitude pairs.
See also
Radio-frequency sweep
Wobbulator
References
Electrical circuits | Sweep generator | [
"Engineering"
] | 518 | [
"Electrical engineering",
"Electronic engineering",
"Electrical circuits"
] |
23,477,549 | https://en.wikipedia.org/wiki/Scout%20X-2M | Scout X-2M was an American expendable launch system which was flown three times between May 1962 and April 1963. It was a four-stage rocket, based on the earlier Scout X-2, but with an MG-18 upper stage instead of the Altair used on the X-2. It was a member of the Scout family of rockets.
The Scout X-2 was an all-solid rocket, with an Algol 1D first stage, a Castor 1A second stage, an Antares 2A third stage, and an MG-18 fourth stage. It was launched from Launch Complex D at Point Arguello, and was used for the launch of P-35 weather satellites.
The first Scout X-2M was launched 24 May 1962, carrying P35-1, but failed to reach orbit. The second flight, launched at 11:44 GMT on 23 August, was the only successful launch to be made by an X-2M, placing P35-2 into low Earth orbit. The final launch, with P35-4 occurred on 26 April 1963, and like the first flight, it failed to reach orbit.
References
1962 in spaceflight
1963 in spaceflight
X-2 | Scout X-2M | [
"Astronomy"
] | 247 | [
"Rocketry stubs",
"Astronomy stubs"
] |
31,019,140 | https://en.wikipedia.org/wiki/Feedthrough | A feedthrough is a conductor used to carry a signal through an enclosure or printed circuit board. Like any conductor, it has a small amount of capacitance. A "feedthrough capacitor" has a guaranteed minimum value of built in it and is used for bypass purposes in ultra-high-frequency applications. Feedthroughs can be divided into power and instrumentation categories. Power feedthroughs are used to carry either high current or high voltage. Instrumentation feedthroughs are used to carry electrical signals (including thermocouples) which are normally low current or voltage. Another special type is what is commonly known as RF-feedthrough, specifically designed to carry very high frequency RF or microwave electrical signals.
A feedthrough electrical connection may have to withstand considerable pressure difference across its length. Systems that operate under high vacuum, such as electron microscopes, require electrical connections through the pressure vessel. Similarly, submersible vehicles require feedthrough connections between exterior instruments and devices and the controls within the vehicle pressure hull. A very common example of a feedthrough connection is an automobile spark plug where the body of the plug must resist the pressure and temperature produced in the engine, while providing a reliable electrical connection to the spark gap in the combustion chamber. (Spark plugs are occasionally used as low-cost or improvised feedthrough connections in non-engine applications.)
There are electrical hermetically sealed feedthroughs for instrumentation, high amperage and voltage, coaxial, thermocouple and fiber optics. Rotary or mechanical feedthroughs also exist.
See also
Vacuum technology
References
Electrical components
de:Leistungstransformator#Durchf.C3.BChrungen | Feedthrough | [
"Technology",
"Engineering"
] | 370 | [
"Electrical engineering",
"Electrical components",
"Components"
] |
31,021,099 | https://en.wikipedia.org/wiki/List%20of%20commercially%20available%20roofing%20materials | Roofing material is the outermost layer on the roof of a building, sometimes self-supporting, but generally supported by an underlying structure. A building's roofing material provides shelter from the natural elements. The outer layer of a roof shows great variation dependent upon availability of material, and the nature of the supporting structure. Those types of roofing material which are commercially available range from natural products such as thatch and slate to commercially produced products such as tiles and polycarbonate sheeting. Roofing materials may be placed on top of a secondary water-resistant material called underlayment.
Steep slope roofing materials
Steep roof materials are roofs that are only recommended where water can freely and openly drain off the edge of the roof without retaining water for too long. The aim is to drain off water completely due to the high water permeability of most of these materials, for instance: the Thatched roofs. In areas where the International Build Code or similar is utilized, the minimum slope required is 2:12, though some countries extend this as high as 4:12
Thatch
Thatch roofing is typically made of plant stalks in overlapping layers.
Wheat straw, widely used in England, France, and other parts of Europe
Seagrass, used in coastal areas where there are estuaries such as Scotland. Has a longer life than straw. Claimed to have a life in excess of 60 years.
Rye straw, commonly used in a barn.
Raffia palm leaves; a well organised raffia palm leaves is mainly used as roof houses in Nigeria, especially among the Igbo.
Rice straw, commonly used in Eastern Asia.
Water reed, commonly used in Ireland for thatching.
Shingle
A shingle is the generic term for an individual roofing unit that is applied with other such units in an overlapping fashion.
Wood shingle, shingles sawn from bolts of wood such as red cedar which has a useful performance life expectancy of up to 30 years. However, young growth red cedar has a short life expectancy and high cost. In the United States and Canada, eastern white cedar is also used. Some hardwoods were very durable roofing found in Colonial Australian and American colonial architecture; their use is now usually limited to building restoration. All wood shingles benefit by being allowed to breathe (dry out from below).
Shake (shingle), Are different than wood shingles in that they are split on one side and sawed on the back side. Commonly referred to as "resawn shakes". A cedar shake is not the same as a cedar shingle.
Asphalt shingle made of bitumen embedded in an organic or fiberglass mat, usually covered with colored, man-made ceramic grit. Cheaper than slate or tiles. The reduced cost of this particular style of roofing is especially apparent in its application and removal. Installation is very streamlined and a rapid process. Depending on the size of the roof and the experience of the crew, it is possible to remove old shingles and apply new ones on 2-3 houses in one day. Life span varies. Use only on slanted roofs.
Rubber shingle, an alternative to asphalt shingle, slate, shake or tile. Made primarily of rubber, often recycled tire-derived rubber. Other typical ingredients include binders, UV (ultraviolet light) inhibitors and color. Warranted and designed to last at least 50 years in most cases.
Asbestos shingles. Very long lifespan, fireproof, and low cost but now rarely used because of health concerns.
Stone slab. Heavy stone slabs (not to be confused with slate) 1–2 inches thick were formerly used as roofing tiles in some regions in England, the Alps, and Scandinavia. Stone slabs require a very heavyweight roof structure, but their weight makes them stormproof. An obsolete roofing material, now used commercially only for building restoration.
Collyweston stone slate named after the village of Collyweston
Solar shingle
Metal shakes or shingles. Long life. High cost, suitable for roofs of 3:12 pitch or greater. Because of the flexibility of metal, they can be manufactured to lock together, giving durability and reducing assembly time. For a discussion of copper system shingles, see Copper in architecture#Wall cladding.
Slate
While slates have high cost, they have a life expectancy of 80 to 400 years. See the article slate industry for an overview including names of quarries. Some of the famous quarries where the highest quality slate comes from that are available in Australia are Bethesda in Wales and areas of Spain.
Ceramic tile
Tile roofing traditionally consists of locally available materials such as clay, granite, terracotta or slate, though many modern applications contain concrete.
Imbrex and tegula, style dating back to ancient Greece and Rome.
Monk and nun, a style similar to Imbrex and tegula, but basically using two Imbrex tiles.
Dutch roof tiles, Netherlands
Mangalore tiles, India
Metal roofing
Metal roofing is any of a large variety of roof coverings made from metal and is characterized by its high resistance, impermeability, and longevity. While there are an infinite variety of how to produce metal roofing, thicknesses, and types for metal/finishes used, roofing is generally grouped into 2 categories: Exposed Fastener Panels and Hidden Fastener Panels. Exposed Fastener panels are held down by fasteners through the outside of the metal, whereas Hidden Fastener Panels are held by hidden fasteners, clips, and sometimes adhesives. Typical metals include Galvanized Steel, Galvalume, Aluminum, Copper, or Vinyl 9Which while not metal is included in many cases for its matching profiles).
Corrugated galvanised iron is galvanised steel manufactured with wavy corrugations to resist lateral flexing and fitted with exposed fasteners. Widely used for low cost and durability. Sheds are normally roofed with this material. Gal iron or Corro was the most extensively used roofing material of 20th century Australia, now replaced in popularity by steel with longer-lasting, coloured, alloy coatings.
Copper roofs can last for hundreds of years. Copper roofing offers durability, ease of fabrication, lighter weight than some other roofing materials, can be curved, low maintenance, corrosion resistance, low thermal movement, lightning protection, radio frequency shielding, and are 100% recyclable. Copper roofs have a high initial cost but very long lifetime: tests on European copper roofs from the 18th Century showed that, in theory, copper roofs can last one thousand years. Another advantage of copper roofing systems is that they are relatively easy to repair.
Standing-seam metal roof with concealed fasteners.
Mechanically seamed metal with concealed fasteners contains sealant in seams for use on very low sloped roofs, suitable for roofs of low pitch such as 0.5/12 to 3/12 pitch.
Flat-seam metal with or without soldered seams.
Steel coated with a coloured alloy of zinc and aluminium.
Stone-coated metal roofing.
Low slope roofing materials
Low slope materials include roofing materials that can be installed at below 2:12 slope, although in the majority of cases low slope materials can also be installed at steeper slopes.
Membrane roofing
Membrane roofing consists of large sheets, generally fused in some way at the joints to form a continuous surface.
Cured Thermoset membrane (e.g. EPDM rubber, Neoprene). Synthetic rubber Cured Thermosets are synthetic rubbers that have undergone the vulcanization or "Curing" process. Seams of materials are bonded by adhesives or chemicals, which over time weaken and separate unless maintained or reinforced. The most commonly used Cured Elastomer membranes are Ethylene Propylene Diene Monomer (commonly EPDM) and Neoprene.
Uncured Thermoset membrane (e.g. CSPE, CPE, NBP, PIB) Uncured elastomers are installed in a manner similar to thermoplastics in that they can be heat or solvent welded. The material then cures over time once exposed to the elements, and then exhibits the same qualities as vulcanized elastomers. The most commonly used Uncured Elastomers are Chlorosulfonated Polyethylene (CSPE), Chlorinated Polyethylene (CPE), Polyisobutylene (PIB), Nitrile Butadiene Polymer (NBP).
Thermoplastics (e.g. PVC, TPO) – Plastic sheets welded together with hot air, creating one continuous sheet membrane. Lends itself well to both big box and small roof application because of its hot air weldability. This membrane is installed by two methods: 1.) Rolls of membrane are attached to the ridged insulation using a bonding adhesive; 2.) The edge of each roll is fastened through ridged insulation into structural deck, and the proceeding roll is lapped over the fasteners. The overlap is then heat-welded with hot air to create a mechanically fastened thermoplastic roof. PVC is also known as IB.Vinyl roof membrane.
Liquid roofing
Liquid roofing includes:
Asphalt roll roofing including single and double coverage types.
Acrylic Based liquid roofing.
Silicone based liquid roofing.
Neoprene Based liquid roofing.
Butyl/Rubber Based liquid roofing.
Modified bitumen
Modified bitumen are long rolls of asphalt-based materials, that can be heat-welded, self-adhering asphalt-adhered, or installed with adhesive. Asphalt is mixed with polymers such as APP or SBS, then applied to fiberglass and/or polyester mat, seams sealed by locally melting the asphalt with heat, hot mopping of asphalt, or adhesive. Lends itself well to most applications.
Built-up roof membrane
Built-up roofs (BURs) consist of multiple plies of bitumen-coated organic felt, polyester felt, or coated fiberglass felts. Three to five plies of felt are laminated to each other and to the substrate with hot asphalt, coal tar pitch, or made-for-purpose cold adhesive. Although the roof membrane can be left bare, it is typically covered with a thick flood coat of the bituminous adhesive and covered with gravel, mineral granules, or a reflective coating, each of which protects the BUR from ultraviolet (UV) light degradation (UV causes evaporation of tar and oxidation of asphalt). Gravel not only provides UV protection, it also helps accommodate sudden temperature changes (thermal shock), protects the surface from hail and mechanical damage, and increases the weight of the roof system to resist wind blow-off.
Fabric
Polyester
PTFE, (synthetic fluoropolymer) embedded in fibreglass
Concrete or fibre cement
Concrete roofing is composed of concrete reinforced with fibers of some sort.
Structural concrete can also be used for flat roof constructions. There are three main categories, precast/prestressed, cast-in-place and shell. There are many types of precast/prestressed concrete roofing. The following are the most common types.
Double tees are the most common products for short spans up to
Hollow-core slabs are used when there is a need for flushed ceiling.
T-beams are similar to double tees but can be used for span ranging from to .
Joists and planks are combination of using prestressed joists with prestressed planks. Keystone-shape joists can be used for up to spans and tee-shape joists can be used for up to span.
Other components
Underlayments
Tar paper and felt paper
Synthetic underlayment
Ice and water shield
Insulations and cover boards
Gyspum Roof Boards
Concrete Roof Boards
Expanded Polystyrene (EPS)
Extruded Polystyrene (XPS)
Polyisocyanurate (ISO)
Wood Fiber Insulation
Fire Sheet
Fiberboard
Fiberglass
Mineral Wool
Spray Foam
Drip edge
Drip edge is a metal installed to the edges of a roof deck, after the roofing material is installed. The metal may be galvanized steel, aluminum, PVC, copper and possibly others.
Gallery
See also
Roof
Domestic roof construction
List of roof shapes
Board roof
Building construction
Building insulation
Building envelope
Grouted roof
Sod roof
Birch-bark roof
Stone roofs such as on a clochán or trullo
Mud roofs such as on beehive house
Sedum a plant used in green roofs
References
Structural engineering
Structural system
Materials
roofing material | List of commercially available roofing materials | [
"Physics",
"Technology",
"Engineering"
] | 2,579 | [
"Structural engineering",
"Building engineering",
"Structural system",
"Construction",
"Materials",
"Architecture lists",
"Civil engineering",
"Roofs",
"Matter",
"Architecture"
] |
31,023,913 | https://en.wikipedia.org/wiki/Lute%20%28material%29 | Lute (from Latin Lutum, meaning mud, clay etc.) was a substance used to seal and affix apparatus employed in chemistry and alchemy, and to protect component vessels against heat damage by fire; it was also used to line furnaces. Lutation was thus the act of "cementing vessels with lute".
In pottery, luting is a technique for joining pieces of unfired leather-hard clay together, using a wet clay slip or slurry as adhesive. The complete object is then fired. Large objects are often built up in this way, for example the figures of the Terracotta Army in ancient China. The edges being joined might be scored or cross-hatched to promote adhesion, but clay and water are the only materials used.
Uses
Lute was commonly used in distillation, which required airtight vessels and connectors to ensure that no vapours were lost; thus it was employed by chemists and alchemists, the latter being known to refer to it as "lutum sapientiae" or the "lute of Wisdom".
The earthen and glass vessels commonly employed in these processes were very vulnerable to cracking, both on heating and on cooling; one way of protecting them was by coating the vessels with lute and allowing it to set. One mixture for this purpose included "fat earth" (terra pinguis), Windsor loam, sand, iron filings or powdered glass, and cow's hair.
Another use for lute was to act as a safety valve, preventing the buildup of vapour pressure from shattering a vessel and possibly causing an explosion. For this purpose, a hole was bored in the flask and covered with luting material of a particular composition, which was kept soft so that excessive buildup of vapour would cause it to come away from the vessel, thus releasing the pressure safely. This process could also be performed manually by the operator removing and reaffixing the lute as required. Lute was also used to effect repairs to cracked glass vessels. In The Alchemist’s Experiment Takes Fire, 1687, one alembic is exploding; the luting used to seal a receiving bottle to another alembic can be seen behind the alchemist's upraised arm.
Lute was frequently applied to the joints between vessels (such as retorts and receivers), making them airtight and preventing vapour from escaping; this was especially important for more penetrating "spiritous" vapours and required a mixture that would set hard - such as a mix of quicklime and either egg white or size etc. However a stronger lute had to be used to confine acid vapours, and for this purpose fat earth and linseed oil were mixed to form "fat lute", which could be rolled into cylinders of convenient size, ready for use. Where the vapour was more "aqueous", and less penetrating, strips of paper affixed with sizing would suffice or "bladder long steeped in water".
Another related use for lute was for lining furnaces, and was described as far back as the 16th century by Georg Agricola in his "De re metallica".
Composition
Fat Lute was made of clay mixed with oil and beaten until it had the consistency of putty. It could be stored in a sealed earthenware vessel, which retained moisture and kept the material pliable. An alchemical writer of the 16th century recommended a lute made up of "loam mixed to a compost with horse dung" while the French chemist Chaptal used a similar mixture of "fat earth" and horse dung, mixed in water and formed into a soft paste.
Linseed meal or Almond meal could be made into a lute by mixing with water or dissolved starch or weak glue, and used in combination with strips of rag or moistened bladder; however, it was combustible which limited its range of applications.
Lime could be made into an effective lute by mixing it with egg white or glue; for sealing joints it was used in conjunction with strips of rag.
Linen rags mixed with paste, or strips of Bladder soaked in warm water, then coated with paste or egg white, also served as a lute.
Fire Lute was used to protect vessels from heat damage. It consisted of clay mixed with sand and either horse-hair or straw or tow (coarse, broken fibre of crops such as flax, hemp, or jute). It had to be allowed to dry thoroughly before use to be effective.
Fusible lute was used to coat earthenware vessels to ensure impermeability. A mixture of Borax and slaked lime, mixed with water into a fine paste, served this purpose.
Parker's Cement, Plaster of Paris and Fusible fluxes (a clay and Borax mixture in 10:1 proportion, mixed to a paste in water) could all be used as lutes, rendering heat protection and air-tightness. Stourbridge clay mixed with water could withstand the highest heat of any lute.
Hard cement was also commonly used to join glass vessels and fix cracks; it was composed of resin, beeswax and either brick dust or "bole earth", or red ochre or venetian red. Soft cement, made of yellow wax, turpentine and venetian red, was also used for repair.
References
Further reading
Antoine-Laurent Lavoisier. Elements of chemistry etc. (Courier Dover Publications, 1790) ch. 7, "Of the composition and application of lutes".
Samuel Frederick Gray & Arthur Livermore Porter. The chemistry of the arts etc. (Carey & Lea, 1830) p. 217 ff., "Chemical lutes".
History of chemistry
Alchemical substances
Distillation
Materials
Joining | Lute (material) | [
"Physics",
"Chemistry"
] | 1,206 | [
"Separation processes",
"Alchemical substances",
"Materials",
"Distillation",
"Matter"
] |
31,025,902 | https://en.wikipedia.org/wiki/Kobayashi%20metric | In mathematics and especially complex geometry, the Kobayashi metric is a pseudometric intrinsically associated to any complex manifold. It was introduced by Shoshichi Kobayashi in 1967. Kobayashi hyperbolic manifolds are an important class of complex manifolds, defined by the property that the Kobayashi pseudometric is a metric. Kobayashi hyperbolicity of a complex manifold X implies that every holomorphic map from the complex line C to X is constant.
Definition
The origins of the concept lie in Schwarz's lemma in complex analysis. Namely, if f is a holomorphic function on the open unit disc D in the complex numbers C such that f(0) = 0 and |f(z)| < 1 for all z in D, then the derivative f '(0) has absolute value at most 1. More generally, for any holomorphic map f from D to itself (not necessarily sending 0 to 0), there is a more complicated upper bound for the derivative of f at any point of D. However, the bound has a simple formulation in terms of the Poincaré metric, which is a complete Riemannian metric on D with curvature −1 (isometric to the hyperbolic plane). Namely: every holomorphic map from D to itself is distance-decreasing with respect to the Poincaré metric on D.
This is the beginning of a strong connection between complex analysis and the geometry of negative curvature. For any complex space X (for example a complex manifold), the Kobayashi pseudometric dX is defined as the largest pseudometric on X such that
,
for all holomorphic maps f from the unit disc D to X, where denotes distance in the Poincaré metric on D. In a sense, this formula generalizes Schwarz's lemma to all complex spaces; but it may be vacuous in the sense that the Kobayashi pseudometric dX may be identically zero. For example, it is identically zero when X is the complex line C. (This occurs because C contains arbitrarily big discs, the images of the holomorphic maps fa: D → C given by f(z) = az for arbitrarily big positive numbers a.)
A complex space X is said to be Kobayashi hyperbolic if the Kobayashi pseudometric dX is a metric, meaning that dX(x,y) > 0 for all x ≠ y in X. Informally, this means that there is a genuine bound on the size of discs mapping holomorphically into X. In these terms, Schwarz's lemma says that the unit disc D is Kobayashi hyperbolic, and more precisely that the Kobayashi metric on D is exactly the Poincaré metric. The theory becomes more interesting as more examples of Kobayashi hyperbolic manifolds are found. (For a Kobayashi hyperbolic manifold X, the Kobayashi metric is a metric intrinsically determined by the complex structure of X; it is not at all clear that this should ever happen. A real manifold of positive dimension never has an intrinsic metric in this sense, because its diffeomorphism group is too big to allow that.)
Examples
Every holomorphic map f: X → Y of complex spaces is distance-decreasing with respect to the Kobayashi pseudometrics of X and Y. It follows that if two points p and q in a complex space Y can be connected by a chain of holomorphic maps C → Y, then dY(p,q) = 0, using that dC is identically zero. This gives many examples of complex manifolds for which the Kobayashi pseudometric is identically zero: the complex projective line CP1 or more generally complex projective space CPn, C−{0} (using the exponential function C → C−{0}), an elliptic curve, or more generally a compact complex torus.
Kobayashi hyperbolicity is preserved under passage to open subsets or to closed complex subspaces. It follows, for example, that any bounded domain in Cn is hyperbolic.
A complex space is Kobayashi hyperbolic if and only if its universal covering space is Kobayashi hyperbolic. This gives many examples of hyperbolic complex curves, since the uniformization theorem shows that most complex curves (also called Riemann surfaces) have universal cover isomorphic to the disc D. In particular, every compact complex curve of genus at least 2 is hyperbolic, as is the complement of 2 or more points in C.
Basic results
For a Kobayashi hyperbolic space X, every holomorphic map C → X is constant, by the distance-decreasing property of the Kobayashi pseudometric. This is often the most important consequence of hyperbolicity. For example, the fact that C minus 2 points is hyperbolic implies Picard's theorem that the image of any nonconstant entire function C → C misses at most one point of C. Nevanlinna theory is a more quantitative descendant of Picard's theorem.
Brody's theorem says that a compact complex space X is Kobayashi hyperbolic if and only if every holomorphic map C → X is constant. An application is that hyperbolicity is an open condition (in the Euclidean topology) for families of compact complex spaces. Mark Green used Brody's argument to characterize hyperbolicity for closed complex subspaces X of a compact complex torus: X is hyperbolic if and only if it contains no translate of a positive-dimensional subtorus.
If a complex manifold X has a Hermitian metric with holomorphic sectional curvature bounded above by a negative constant, then X is Kobayashi hyperbolic. In dimension 1, this is called the Ahlfors–Schwarz lemma.
The Green–Griffiths–Lang conjecture
The results above give a complete description of which complex manifolds are Kobayashi hyperbolic in complex dimension 1. The picture is less clear in higher dimensions. A central open problem is the Green–Griffiths–Lang conjecture: if X is a complex projective variety of general type, then there should be a closed algebraic subset Y not equal to X such that every nonconstant holomorphic map C → X maps into Y.
Clemens and Voisin showed that for n at least 2, a very general hypersurface X in CPn+1 of degree d at least 2n+1 has the property that every closed subvariety of X is of general type. ("Very general" means that the property holds for all hypersurfaces of degree d outside a countable union of lower-dimensional algebraic subsets of the projective space of all such hypersurfaces.) As a result, the Green–Griffiths–Lang conjecture would imply that a very general hypersurface of degree at least 2n+1 is Kobayashi hyperbolic. Note that one cannot expect all smooth hypersurfaces of a given degree to be hyperbolic, for example because some hypersurfaces contain lines (isomorphic to CP1). Such examples show the need for the subset Y in the Green–Griffiths–Lang conjecture.
The conjecture on hyperbolicity is known for hypersurfaces of high enough degree, thanks to a series of advances by Siu, Demailly and others, using the technique of jet differentials. For example, Diverio, Merker and Rousseau showed that a general hypersurface in CPn+1 of degree at least 2n5 satisfies the Green-Griffiths-Lang conjecture. ("General" means that this holds for all hypersurfaces of given degree outside a finite union of lower-dimensional algebraic subsets of the projective space of all such hypersurfaces.) In 2016, Brotbek gave a proof of the Kobayashi conjecture for the hyperbolicity of general hypersurfaces of high degree, based on a use of Wronskian differential equations; explicit degree bounds have then been obtained in arbitrary dimension by Ya Deng and Demailly, e.g. [(en)2n+2/3] by the latter. Better bounds for the degree are known in low dimensions.
McQuillan proved the Green–Griffiths–Lang conjecture for every complex projective surface of general type whose Chern numbers satisfy c12 > c2. For an arbitrary variety X of general type, Demailly showed that every holomorphic map C→ X satisfies some (in fact, many) algebraic differential equations.
In the opposite direction, Kobayashi conjectured that the Kobayashi pseudometric is identically zero for Calabi–Yau manifolds. This is true in the case of K3 surfaces, using that every projective K3 surface is covered by a family of elliptic curves. More generally, Campana gave a precise conjecture about which complex projective varieties X have Kobayashi pseudometric equal to zero. Namely, this should be equivalent to X being special in the sense that X has no rational fibration over a positive-dimensional orbifold of general type.
Analogy with number theory
For a projective variety X, the study of holomorphic maps C → X has some analogy with the study
of rational points of X, a central topic of number theory. There are several conjectures on the relation between these two subjects. In particular, let X be a projective variety over a number field k. Fix an embedding of k into C. Then Lang conjectured that the complex manifold X(C) is Kobayashi hyperbolic if and only if X has only finitely many F-rational points for every finite extension field F of k. This is consistent with the known results on rational points, notably Faltings's theorem on subvarieties of abelian varieties.
More precisely, let X be a projective variety of general type over a number field k. Let the exceptional set Y be the Zariski closure of the union of the images of all nonconstant holomorphic maps C → X. According to the Green–Griffiths–Lang conjecture, Y should be a proper closed subset of X (and, in particular, not be equal to X). The strong Lang conjecture predicts that Y is defined over k and that X − Y has only finitely many F-rational points for every finite extension field F of k.
In the same spirit, for a projective variety X over a number field k (or, more generally, a finitely generated field k of characteristic zero), Campana conjectured that the Kobayashi pseudometric of X(C) is identically zero if and only if X has potentially dense rational points, meaning that there is a finite extension field F of k such that the set X(F) of F-rational points is Zariski dense in X.
Variants
The Carathéodory metric is another intrinsic pseudometric on complex manifolds, based on holomorphic maps to the unit disc rather than from the unit disc. The Kobayashi infinitesimal pseudometric is a Finsler pseudometric whose associated distance function is the Kobayashi pseudometric as defined above. The Kobayashi–Eisenman pseudo-volume form is an intrinsic measure on a complex n-fold, based on holomorphic maps from the n-dimensional polydisc to X. It is understood better than the Kobayashi pseudometric. In particular, every projective variety of general type is measure-hyperbolic, meaning that the Kobayashi–Eisenman pseudo-volume form is positive outside a lower-dimensional algebraic subset.
Analogous pseudometrics have been considered for flat affine and projective structures, as well as for more general projective connections and conformal connections.
Notes
References
Algebraic geometry
Complex manifolds | Kobayashi metric | [
"Mathematics"
] | 2,364 | [
"Fields of abstract algebra",
"Algebraic geometry"
] |
39,117,935 | https://en.wikipedia.org/wiki/Mechanical%20properties%20of%20biomaterials | Materials that are used for biomedical or clinical applications are known as biomaterials. The following article deals with fifth generation biomaterials that are used for bone structure replacement. For any material to be classified for biomedical applications, three requirements must be met. The first requirement is that the material must be biocompatible; it means that the organism should not treat it as a foreign object. Secondly, the material should be biodegradable (for in-graft only); the material should harmlessly degrade or dissolve in the body of the organism to allow it to resume natural functioning. Thirdly, the material should be mechanically sound; for the replacement of load-bearing structures, the material should possess equivalent or greater mechanical stability to ensure high reliability of the graft.
Introduction
The biomaterial term is used for materials that can be used in biomedical and clinical applications. They are bioactive and biocompatible in nature. Currently, many types of metals and alloys (stainless steel, titanium, nickel, magnesium, Co–Cr alloys, Ti alloys), ceramics (zirconia, bioglass, alumina, hydroxyapatite) and polymers (acrylic, nylon, silicone, polyurethane, polycaprolactone, polyanhydrides) are used for load bearing applications. This includes dental replacements and bone joining or replacements for medical and clinical application. Therefore, their mechanical properties are very important. Mechanical properties of some biomaterials and bone are summarized in Table 1. Among them, hydroxyapatite is most widely studied bioactive and biocompatible material. However, it has lower Young's modulus and fracture toughness with a brittle nature. Hence, it is required to produce a biomaterial with good mechanical properties.
Elastic modulus
Elastic modulus is simply defined as the ratio of stress to strain within the proportional limit. Physically, it represents the stiffness of a material within the elastic range when tensile or compressive loads are applied. It is clinically important because it indicates the selected biomaterial has similar deformable properties with the material it is going to replace. These force-bearing materials require high elastic modulus with low deflection. As the elastic modulus of material increases, fracture resistance decreases. It is desirable that the biomaterial elastic modulus is similar to that of bone. This is because if it is more than bone's elastic modulus then the load is borne by the material only; while the load is borne by bone only if it is less than bone material. The elastic modulus of a material is generally calculated by the bending test, because deflection can be easily measured in this case as compared to very small elongation in compressive or tensile load. However, biomaterials (for bone replacement) are usually porous and the sizes of the samples are small. Therefore, the nanoindentation test is used to determine the elastic modulus of these materials. This method has high precision and is convenient for micro-scale samples. Another method of elastic modulus measurement is the non-destructive method. It is also a clinically very good method because of its simplicity and repeatability since materials are not destroyed.
Hardness
Hardness is a measure of plastic deformation and is defined as the force per unit area of indentation or penetration. Hardness is one of the most important parameters for comparing properties of materials. It is used for finding the suitability of the clinical use of biomaterials. Biomaterial hardness is desirable as equal to bone hardness. If higher than the biomaterial, then it penetrates in the bone. Higher hardness results in less abrasion. As said above, biomaterials sample are very small, therefore micro- and nano-scale hardness tests (Diamond Knoop and Vickers indenters) are used.
Fracture strength
The strength of a material is defined as the maximum stress that can be endured before fracture occurs. Strength of biomaterials (bioceramics) is an important mechanical property because they are brittle. In brittle materials like bioceramics, cracks easily propagate when the material is subject to tensile loading, unlike compressive loading. A number of methods are available for determining the tensile strength of materials, such as the bending flexural test, the biaxial flexural strength test and the weibull approach. In bioceramics, flaws influence the reliability and strength of the material during implantation and fabrication. There are a number of ways that flaws can be produced in bioceramics such as thermal sintering and heating. It is important for bioceramics to have high reliability, rather than high strength.
The strength of brittle materials depends on the size of flaws distributed throughout the material. According to Griffith's theory of fracture in tension, the largest flaw or crack will contribute the most to the failure of a material. Strength also depends on the volume of a specimen since flaw size is limited to the size of the specimen's cross section. Therefore, the smaller the specimen (e.g., fibers), the higher the fracture strength. Porosity of implanted bioceramic has a tremendous influence on the physical properties. Pores are usually formed during processing of materials. Increasing the porosity and pore size means increasing the relative void volume and decreasing density; this leads to a reduction in mechanical properties and lowers the overall strength of bioceramic.
To use ceramics as self-standing implants that are able to withstand tensile stresses is a primary engineering design objective. Four general approaches have been used to achieve this objective:
1) use of the bioactive ceramic as a coating on a metal or ceramic substrate
2)strengthening of the ceramic, such as via crystallization of glass
3) use of fracture mechanics as a design approach and
4) reinforcing of the ceramic with a second phase.
For example, hydroxyapatite and other calcium phosphates bioceramics are important for hard tissue repair because of their similarity to the minerals in natural bone, and their excellent biocompatibility and bioactivity, but they have poor fatigue resistance and strength. Hence, bioinert ceramic oxides with high strength are used to enhance the densification and the mechanical properties of them.
Fracture toughness
Fracture toughness is required to alter the crack propagation in ceramics. It is helpful to evaluate the serviceability, performance and long term clinical success of biomaterials. It is reported that the high fracture toughness material improved clinical performance and reliability as compare to low fracture toughness. It can be measured by many methods e.g. indentation fracture, indentation strength, single edge notched beam, single edge pre cracked beam and double cantilever beam.
Fatigue
Fatigue is defined as failure of a material due to repeated/cyclic loading or unloading (tensile or compressive stresses). It is also an important parameter for biomaterial because cyclic load is applied during their serving life. In this cyclic loading condition, micro crack/flaws may be generated at the interface of the matrix and the filler. This micro crack can initiate permanent plastic deformation which results in large crack propagation or failure. During the cyclic load several factor also contribute to microcrack generation such as frictional sliding of the mating surface, progressive wear, residual stresses at grain boundaries, stress due to shear.
Table 1: Summary of mechanical properties of cortical bone and biomaterial
Fatigue fracture and wear have been identified as some of the major problems associated with implant loosening, stress-shielding and ultimate implant failure. Although wear is commonly reported in orthopaedic applications such as knee and hip joint prostheses, it is also a serious and often fatal experience in mechanical heart valves.
The selection of biomaterials for wear resistance unfortunately cannot rely only on conventional thinking of using hard ceramics, because of their low coefficient of friction and high modulus of elasticity. This is because ceramics are generally prone to brittle fracture (having a fracture toughness typically less than 1 MPa√m) and need absolute quality control to avoid fatigue fracture for medical device applications. The development of fatigue fracture and wear resistant biomaterials looks into the biocomposites of two or more different phases such as in interpenetrating network composites. The advantage of these composites is that one can incorporate controlled drug release chemicals, friction modifiers, different morphologies to enable better host–implant performance and chemical entities to reduce or aid removal of wear debris. Of equal importance are the tools developed to predict fatigue fracture/wear using new methodologies involving in vitro tests, computational modelling to obtain design stresses and fracture/wear maps to identify mechanisms.
Viscoelasticity
Viscoelasticity, a material property characterized by the extrusion of dual solid and liquid-like behaviors, is typically found in an array of polymer-based biomaterials, including those used in biomedical devices as well as in clinical settings. From polymer-based surface coatings on drug-eluting stents to entangled tissue networks that have load-bearing capabilities and hydrogels that possess complex crosslinks, all of these examples display viscoelastic behavior. Often times, flow plasticity theory and linear elasticity are utilized to describe the rheological behavior of metals and other hard materials, yet they are not commonly used to elaborate on the material behavior of biomaterials. Viscoelasticity is often described in terms of its time-dependent material properties associated with its characteristic stress relaxation time. Additionally, the energy dissipation associated with the liquid-like portion of the response to an applied load can be funneled into the complex modulus, which is represented by two distinct categories: one real and one imaginary, for the viscoelastic response[10]. The viscoelastic response of a biomaterial can be modeled by linear mathematical models, and atypically a non-linear mathematical models that corresponds to the loading capabilities of the biomaterial in use.
Viscoelasticity in polymeric biomaterials
There is a tendency for polymeric biomaterials to display the same characteristics as solid, rigid materials over a short time span, in addition to exhibiting exceptional flow behavior over longer periods of time. This translates to long-term analysis and studies focused towards ensuring the mechanical integrity of these biomaterials to prevent potential deformation and mechanical failure once employed in a clinical setting. The viscoelastic behavior is typically dependent on factors such as the crosslink density, the average molecular weight, the degree of crystallinity, and the degree of entanglement as well as the general chemistry of the biomaterial. There are modeling programs employed to probe the material behavior over an array of temperatures and applied frequencies, as well as to decrease the potential for complexity in synthesizing polymers at the industrial level and for commercial use. The programs themselves often focus on decreasing the rate of mechanical and environmental degradation by focusing on probing the rate sensitivity as well as creep response[11]. For example, in polymeric grafts that act as replacements for tissues, the viscoelastic response is necessary to be mimicked to ensure ample biocompatibility and structural stability over the life-span of the material.
Viscoelasticity in tissue
Tissues themselves are, at their fundamental level, an amalgamation of entangled and crosslinked polymer networks that are composed of collagen, other organic compounds found in the human body, and long polymer chain structures. The degree at which entanglement occurs, crosslinking behavior between other compounds, and the interpenetration ability of excess polymer networks determines the outstanding character of a tissue network. Everything from macroscopic structural to atomistic-level arrangements in a tissue, such as the crimping behavior seen in tendons, can give way to nonlinear elastic behavior which can be highly expressed due to the intermolecular arrangements within the material. Since tissues are hydrolyzed to maintain biological function, this often affects their mechanical performance as it often results in the liquid component affecting the deformation response of the material. Additionally, the degree of crosslinking present in an individual crosslinked collagen network can be prone to the biological environment of said crosslinked network[12]. With that in mind, the time-dependent mechanical properties of tissues can be incredibly interdependent on molecular interactions and the chemical environment in which a specific tissue is native to. In comparison to other tissue, articular cartilage itself begins to enlarge when subjected to unloading and this puts the microstructure of the material into a state of tension. Articular cartilage, a native biomaterial, typically supplies a soft base for tail end of narrow bones located in synovial joints while providing lubrication capabilities that allow joints to interact without excess friction. The cartilage itself is composed of collagen fiber within an entangled gel-like structure. This tissue structure behaves similar to a viscoelastic solid in the sense that the response to strain under an excess load is dependent on the rate of the load. Furthermore, when a mechanical load is applied to the tissue, the fluid is forced out of the porous membranes of the biomaterial which exacerbates permanent deformation, while simultaneously stifling viscous flow and decreasing energy in the material overall. Overall, the viscoelastic characteristics and the viscous attributes in the liquid phase play a role in the dynamic behavior of tissue, and tissue-based materials.
See also
Artificial bone
Biomaterials
Stress
Strain
Hooke's law
Shear modulus
Bending stiffness
Toughness
References
Further reading
Ward, I.M. (1983). The Mechanical Properties of Solid Polymers. New York: Wiley.
Sychterz, C.J., Yang, A., and Engh, C.A. (1999). "Analysis of temporal wear patterns of porous-coated acetabular components: Distinguishing between true wear and so-called bedding-in". Journal of Bone and Joint Surgery (American), 81A(6), 821–30.
Saito, M., and Marumo, Y. (2010). "Collagen cross-links as a determinant of bone quality: a possible explanation for bone fragility in aging, osteoporosis, and diabetes mellitus". Osteoporosis International 21(2), 195–214.
Ichim, Q. Li, W. Li, M.V. Swain, J. Kieser (2007). "Modelling of fracture behaviour in biomaterials". Biomaterials 28(7). 1317–1326.
S.H Teoh (2000). "Fatigue of biomaterials: a review". International Journal of Fatigue 22(10). 825–837.
Bhatia, S. K. (2010). Biomaterials for clinical applications. Springer.
Hench, L. L. (1993). An introduction to bioceramics (Vol. 1). World Scientific.
Biomaterials | Mechanical properties of biomaterials | [
"Physics",
"Biology"
] | 3,096 | [
"Biomaterials",
"Materials",
"Matter",
"Medical technology"
] |
39,124,477 | https://en.wikipedia.org/wiki/Combined%20forced%20and%20natural%20convection | In fluid thermodynamics, combined forced convection and natural convection, or mixed convection, occurs when natural convection and forced convection mechanisms act together to transfer heat. This is also defined as situations where both pressure forces and buoyant forces interact. How much each form of convection contributes to the heat transfer is largely determined by the flow, temperature, geometry, and orientation. The nature of the fluid is also influential, since the Grashof number increases in a fluid as temperature increases, but is maximized at some point for a gas.
Characterization
Mixed convection problems are characterized by the Grashof number (for the natural convection) and the Reynolds number (for the forced convection). The relative effect of buoyancy on mixed convection can be expressed through the Richardson number:
The respective length scales for each dimensionless number must be chosen depending on the problem, e.g. a vertical length for the Grashof number and a horizontal scale for the Reynolds number. Small Richardson numbers characterize a flow dominated by forced convection. Richardson numbers higher than indicate that the flow problem is pure natural convection and the influence of forced convection can be neglected.
Like for natural convection, the nature of a mixed convection flow is highly dependent on heat transfer (as buoyancy is one of the driving mechanisms) and turbulence effects play a significant role.
Cases
Because of the wide range of variables, hundreds of papers have been published for experiments involving various types of fluids and geometries. This variety makes a comprehensive correlation difficult to obtain, and when it is, it is usually for very limited cases. Combined forced and natural convection, however, can be generally described in one of three ways.
Two-dimensional mixed convection with aiding flow
The first case is when natural convection aids forced convection. This is seen when the buoyant motion is in the same direction as the forced motion, thus accelerating the boundary layer and enhancing the heat transfer. Transition to turbulence, however, can be delayed. An example of this would be a fan blowing upward on a hot plate. Since heat naturally rises, the air being forced upward over the plate adds to the heat transfer.
Two-dimensional mixed convection with opposing flow
The second case is when natural convection acts in the opposite way of the forced convection. Consider a fan forcing air upward over a cold plate. In this case, the buoyant force of the cold air naturally causes it to fall, but the air being forced upward opposes this natural motion. Depending on the Richardson number, the boundary layer at the cold plate exhibits a lower velocity than the free stream, or even accelerates in the opposite direction. This second mixed convection case therefore experiences strong shear in the boundary layer and quickly transitions into a turbulent flow state.
Three-dimensional mixed convection
The third case is referred to as three-dimensional mixed convection. This flow occurs when the buoyant motion acts perpendicular to the forced motion. An example of this case is a hot, vertical flate plate with a horizontal flow, e.g. the surface of a solar thermal central receiver. While the free stream continues its motion along the imposed direction, the boundary layer at the plate accelerates in the upward direction. In this flow case, buoyancy plays a major role in the laminar-turbulent transition, while the imposed velocity can suppress turbulence (laminarization)
Calculation of total heat transfer
Simply adding or subtracting the heat transfer coefficients for forced and natural convection will yield inaccurate results for mixed convection. Also, as the influence of buoyancy on the heat transfer sometimes even exceeds the influence of the free stream, mixed convection should not be treated as pure forced convection. Consequently, problem-specific correlations are required. Experimental data has suggested that
can describe the area-averaged heat transfer. For the case of a large, vertical surface in a horizontal flow provided a best fit depending on the details of how is fitted.
Applications
Combined forced and natural convection is often seen in very-high-power-output devices where the forced convection is not enough to dissipate all of the heat necessary. At this point, combining natural convection with forced convection will often deliver the desired results. Examples of these processes are nuclear reactor technology and some aspects of electronic cooling.
References
Convection
Heat transfer | Combined forced and natural convection | [
"Physics",
"Chemistry"
] | 856 | [
"Transport phenomena",
"Physical phenomena",
"Heat transfer",
"Convection",
"Thermodynamics"
] |
39,124,671 | https://en.wikipedia.org/wiki/NK-92 | The NK-92 cell line is an immortalised cell line that has the characteristics of a type of immune cell found in human blood called ’natural killer’ (NK) cells. Blood NK cells and NK-92 cells recognize and attack cancer cells as well as cells that have been infected with a virus, bacteria, or fungus. NK-92 cells were first isolated in 1992 in the laboratory of Hans Klingemann at the British Columbia Cancer Agency in Vancouver, Canada, from a patient who had a rare NK cell non-Hodgkin-lymphoma. These cells were subsequently developed into a continuously growing cell line. NK-92 cells are distinguished by their suitability for expansion to large numbers, ability to consistently kill cancer cells and testing in clinical trials. When NK-92 cells recognize a cancerous or infected cell, they secrete perforin that opens holes into the diseased cells and releases granzymes that kill the target cells. NK-92 cells are also capable of producing cytokines such as tumor necrosis factor alpha (TNF-a) and interferon gamma (IFN-y), which stimulates proliferation and activation of other immune cells.
In clinical trials
Several phase 1 clinical trials have been performed by experts in the field of adoptive immunotherapy of cancer. Hans Klingemann and Sally Arai completed a US trial at Rush University Medical Center (Chicago) in renal cell cancer and melanoma patients in 2008, and Torsten Tonn, MD and Oliver Ottmann, MD completed the European trial at the University of Frankfurt in patients with various solid and hematological malignancies in 2013. Armand Keating at Princess Margaret Hospital in Toronto conducted a trial in which NK-92 cells were given to patients who had relapsed after autologous bone marrow transplants for leukemia or lymphoma. In all clinical trials so far, NK-92 cells were administered as a simple intravenous infusion, dosed two or three times per treatment course, and given in the outpatient setting.
Of the 39 patients enrolled across the three studies, 2 serious (grade 3–4) side-effects occurred during or after the infusion of NK-92 cells, the side effects disappeared afterward. The doses given to patients ranged from 1x108 cells/m2 to 1x1010 cells/m2 per infusion. Patients received between two and three infusions over a period of less than a week. About one-third of the treated patients had clinically meaningful responses with some of them fully recovering.
Comparison to other NK cells
In a 2017 study by Congcong Zhang and Winfried S. Wels, NK-92 cells were genetically engineered to recognize and kill specific human cancers by expressing chimeric antigen receptors (CARs). CAR-engineered T-lymphocytes (CAR-T) have garnered attention in immuno-oncology, as the infusion of CAR-T cells has been shown to induce remissions in some patients with acute and chronic leukemia and lymphoma. However, CAR-T cells can cause cytokine release syndrome (CRS). CAR-engineered NK cells from either peripheral or cord blood have not proved to be as feasible for use to treat diseases as they are difficult to expand to get sufficient numbers, and the yields can be variable and/or too low. Also, genetic transduction to introduce the CAR into blood NK cells requires lentiviral or retroviral vectors, which are only moderately efficient.
NK-92 cells, in contrast to NK-92 CAR-T cells, have predictable expansion kinetics and can be grown in bioreactors that produce billions of cells within a couple of weeks. Further, NK-92 cells can easily be transduced by physical methods, and mRNA can be shuttled into NK-92 cells with high efficiency. CAR-expressing NK-92 have been generated to target a number of cancer surface receptors such as programmed death domain ligand 1 (PD-L1), CD19 (a type of B cell receptor), human epidermal growth factor receptor 2 (HER2/ErbB2) and epidermal growth factor receptor (EGFR, aka HER1); and many of these engineered NK-92 cells are currently in clinical trials for the treatment of cancer.
NK-92 variants
NK-92 cells, which require interleukin-2 (IL-2) for growth, have also been genetically altered with an IL-2 gene to allow them to grow in culture without the addition of IL-2. They have also been engineered to express a high-affinity Fc-receptor which is the main receptor for monoclonal antibodies to bind to NK-92 and use their cytotoxic load to kill cancer cells. The cells have been further engineered to express Chimeric Antigen Receptors (CARs) such as programmed death domain ligand 1 (PD-L1). During the course of development, NK-92 cells were renamed activated NK cells (aNK) and the different variants have been designated as follows:
NK-92 = parental cells, later designated aNK
NK-92ci = NK-92 cells transfected with an episomal vector for expression of IL-2
NK-92 mi = NK-92 cells transfected with an MFG vector for expression of IL-2
haNK = NK-92 (aNK) transfected with a plasmid expressing high affinity CD16 FcR and erIL-2
taNK = NK-92 (aNK) transfected with either a plasmid or lentiviral vector expressing a CAR
t-haNK = NK-92 (aNK) transfected with a plasmid expressing a CAR and CD16 FcR erIL-2
qt-haNK = NK-92 (aNK) transfected with a plasmid expressing a 4th gene in addition to a CAR, the CD16 FcR, and erIL-2: examples: homing receptor of the CXCR family or immune-active cytokines
The high affinity Fc-receptor-expressing NK (haNK) cells were administered to patients with advanced Merkel cell carcinoma (MCC) and there were some notable responses. Currently, a HER2-targeted aNK (taNK) line and various t-haNK (CAR and Fc-receptor expressing) cell lines are in clinical trials in patients with various cancers, as described in the review “The NK-92 cell line 30 years later: its impact on natural killer cell research and treatment of cancer."
Ownership and Licenses
Global rights to the NK-92 cell line were assigned to ImmunityBio Inc. (formerly NantKwest, Inc.). ImmunityBio's only authorized NK-92 distributor is Brink Biologics, Inc. (San Diego), which makes NK-92 cells and certain genetically modified CD16+ variants available to third parties for non-clinical research under a limited use license agreement.
References
External links
Cellosaurus entry for NK-92
Human cell lines
Cancer treatments
Biotechnology
Genetic engineering | NK-92 | [
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41,842,064 | https://en.wikipedia.org/wiki/Tetranitratoaluminate | Tetranitratoaluminate is an anion of aluminium and nitrate groups with formula [Al(NO3)4]− that can form salts called tetranitratoaluminates. It is unusual in being a nitrate complex of a light element.
Related substances
By substituting boron for aluminium tetranitratoborates result. Aluminium can coordinate more nitrates resulting in pentanitratoaluminates and hexanitratoaluminates.
By replacing nitrate with perchlorate, the tetraperchloratoaluminate ion results.
Formation
When hydrated aluminium nitrate reacts with dinitrogen pentoxide it forms a nitronium salt: [NO2]+[Al(NO3)4]−.
A way to make a tetranitratoaluminate salt of a cation is to treat the chloride of the cation and aluminium chloride with liquid dinitrogen tetroxide pure or dissolved in nitromethane. The reaction is started at liquid nitrogen temperatures and then warmed up. Dark red nitrosyl chloride is formed as a byproduct. The byproducts and solvents can then be evaporated. The tetramethylammonium salt can form this way.
Properties
The tetranitratoaluminate group has two bidentate nitrate groups attached in a square around the aluminium, and with two other monodentate nitrates attached via one oxygen only, perpendicular, up and down from the square.
Tetranitratoaluminate salts are not completely stable and can decompose to the nitrates and aluminium oxynitrates.
When nitronium tetranitratoaluminate is sublimed it can form anhydrous aluminium nitrate.
Nitronium tetranitratoaluminate dissolved in a nitric acid and dinitrogen pentoxide mixture yields the hexanitratoaluminate complex. In water is it converted to the hexaaqua complex with six water molecules replacing the nitrate groups.
Examples
Tetraethyl ammonium tetranitratoaluminate along with nitronium tetranitratoaluminate were the first to be discovered.
1-Ethyl-4,5-dimethyl-tetrazolium tetranitratoaluminate is an oxygen balanced ionic liquid, This liquid salt is stable when moisture is excluded. It is soluble in methyl nitrate. It solidifies to a glass at −46 °C, starts to slowly decompose at 75 °C, and inflames, without oxygen required, around 200 °C. When it burns it produces aluminium oxide, nitrogen, water and carbon monoxide. It is being proposed as a rocket propellant, because it has better performance than hydrazine. Rubidium and caesium also form salts.
Tetramethyl ammonium tetranitratoaluminate forms monoclinic crystals with a = 12.195 Å, b = 9.639 Å, c = 12.908 Å, α = 90°, β = 110.41°, γ = 90°, and a formula weight of 349.17 (4 equivalents) per unit cell. The unit cell volume is 1422 Å3 and the calculated density 1.631 g/cm3.
References
Aluminium complexes
Nitrates
Anions | Tetranitratoaluminate | [
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41,848,782 | https://en.wikipedia.org/wiki/Kodaikanal%20mercury%20poisoning | Kodaikanal mercury poisoning is a proven case of mercury contamination at the hill station of Kodaikanal, Tamil Nadu, India by Hindustan Unilever in the process of making mercury thermometers for export around the world. The exposé of the environmental abuse led to the closure of the factory in 2001 and opened up a series of issues in India such as corporate liability, corporate accountability and corporate negligence.
Mercury pollution in Kodaikanal
The mercury contamination in Kodaikanal originated at a thermometer factory owned by Hindustan Unilever. Unilever acquired the thermometer factory from cosmetics maker Pond's India Ltd. Pond's moved the factory from the United States to India in 1982 after the plant owned there by its parent, Chesebrough-Pond's, had to be dismantled following increased awareness in developed countries of polluting industries. In 1987, Pond's India and the thermometer factory went to Hindustan Unilever when it acquired Cheseborough-Pond's globally.
The factory imported mercury from the United States, and exported finished thermometers to markets in the United States and Europe. Around 2001, a number of workers at the factory began complaining of kidney and related ailments. Public interest groups such as Tamil Nadu Alliance Against Mercury (TNAAC) alleged that the Company had been disposing mercury waste without following proper protocols. In early 2001, public interest groups unearthed a pile of broken glass thermometers with remains of Mercury from an interior of part of the shola forest, which they suspected could have come from the company. In March, a public protest led by local workers' union and international environmental organisation Greenpeace forced the company to shut down the factory. Soon the company admitted that it did dispose of mercury contaminated waste. The company said in its 2002 annual report and its latest Sustainability Report that it did not dump glass waste contaminated with mercury on the land behind its factory, but only a quantity of 5.3 metric tonnes of glass containing 0.15% residual mercury had been sold to a scrap recycler located about three kilometers from the factory, in breach of the company procedures. Quoting a report prepared by an international environmental consultant, Unilever said there was no health effect on the workers of the factory or any impact on the environment. This is hotly contested by a book published by Pan MacMillan in 2023, Heavy Metal: How a Global Corporation Poisoned Kodaikanal, authored by veteran journalist-tuned-public policy leader Ameer Shahul.
Reverse dumping
Once the factory was shut down, public interest groups demanded the return of the remaining mercury waste to the United States for recycling, remediation of the factory site, and address of the health complaints of the workers. Local groups and workers' union under the leadership of Greenpeace, represented to the company, regulatory bodies, and the government, besides initiating legal action against the company.
Greenpeace campaigner Ameer Shahul led the public affairs groups and workers collaboration in forcing the Company to collect 290 tonnes of dumped mercury waste from the shola forest and send back to the United States for recycling in 2003. This was widely hailed by the media as ‘reverse dumping’. Later Greenpeace campaigners Ameer Shahul and Navroz Mody led the groups in lobbying for remediation of the site and initiated an investigation by the Department of Atomic Energy of Government of India, which found that the free mercury level in the atmosphere of Kodaikanal was 1000 times more than what is found in normal conditions. Analysis of water, sediment and fish samples collected from Kodaikanal Lake by a team of scientists of the Department of Atomic Energy showed elevated levels of mercury four years after the stoppage of mercury emissions. A series of scientific studies have also been carried out by Governmental and non-governmental organisations to determine the extent of damage caused to the environment and to the people who were exposed to mercury in the factory.
Remediation of the site
Greenpeace and workers' unions continued to mount pressure on the company to take responsibility for the dumping crimes it had committed and for meddling with a pristine environment. They asked the regulatory bodies to prosecute the company. With these demands, public interest groups led by Greenpeace campaign head Shahul spooked the annual general body meeting of Hindustan Unilever in 2004. Consequently, the company began working with the regulatory body Tamil Nadu Pollution Control Board (TNPCB) to remediate the soil, de-contaminate and scrap the thermometer-making equipment at the Kodaikanal site. The company appointed National Environmental Engineering Research Institute (NEERI) to finalise the scope for remediation, which was vehemently opposed by environmentalists. In 2006, the plant, machinery and materials used in thermometer manufacturing at the site were decontaminated and disposed of as scrap to industrial recyclers. In the following year, NEERI conducted trials at the factory for remediation of the contaminated soil on site, and recommended a remediation protocol of soil washing and thermal retorting. These were hotly contested by environmental groups under the leadership of Nityanand Jayaraman. Ultimately, the Tamil Nadu Pollution Control Board (TNPCB) recommended a remediation standard of up to 20 mg/kg of mercury concentration in soil, which means 95% of the samples analysed after the remediation process should be of less than 20 mg/kg. Consequently, pre-remediation work started in May 2009.
Public interest groups contested the soil clean-up criteria and alleged that TNPCB is helping Unilever clean up to lower standards to cut costs. The acceptable mercury level being suggested by TNPCB is at least 20 times higher than what Unilever would have been required to do if they had caused the same contamination in the United Kingdom, where they are based. They also called for transparency and public participation in the process of deciding the levels of clean-up and in the process of clean-up.
Workers' health problems
After the shut down of the factory, the health specialists from Bangalore-based Community Health Centre conducted a survey among the former workers of the factory. It found that former workers of the factory had visible signs of mercury poisoning such as gum and skin allergy and related problems, 'which appeared to be due to exposure to mercury'.
The company claims that comprehensive occupational safety and health systems existed at the Kodaikanal factory prior to its closure in 2001. Internal monitoring within the factory and external audits carried out by statutory authorities during the operations of the factory showed that there were no adverse health effects to the workers on account of their employment at the factory. It says there had been a comprehensive medical examination conducted by a panel of doctors using a questionnaire developed by Mine Safety and Health Administration (MSHA) of the United States Department of Labor; a study by the Certifying Surgeon from the Inspectorate of Factories; an assessment by P N Viswanathan of Indian Institute of Toxicology Research (IITR); a study by Tom van Teuenbroek of TNO; and a study by IITR, formerly known as Industrial Toxicology Research Centre(ITRC) as directed by a Monitoring Committee set up by the Supreme Court of India.
The company says its conclusions of its occupational health surveillance were also endorsed by the All India Institute of Medical Sciences (AIIMS) and the National Institute of Occupational Health (NIOH).
In February 2006, a group of ex-employees of the factory approached the Madras High Court seeking directions for conducting a fresh health survey and providing economic rehabilitation. A year later, the Madras High Court constituted a five-member expert committee, with representatives from ITRC, AIIMS and NIOH to decide whether the alleged health conditions of the workers and their families were related to mercury exposure, and recommend whether there was need for a new health study. The Committee after examining the ex-workers, questioning the Hindustan Unilever Limited (HUL) officials and after a visit to the factory in October 2007 submitted its report suggesting that there is "no sufficient evidence to link the current clinical condition of the factory workers to the mercury exposure in the factory in the past". Accepting the report, the Madras High Court ruled out the need for any fresh health study. In the meantime, the Ministry of Labour and Employment, which is also a respondent in the case before the Madras High Court conducted a detailed study by a team comprising experts from various fields found that there is prima facie evidence to suggest that not only the workers of the factory, but even the children of the workers, have suffered because of exposure to mercury. The Ministry submitted its report to the Madras High Court in 2011. It also recommended setting up a Board to examine the extent of damage or disability suffered by workers and their children because of exposure to mercury, and based on the assessment of the Board workers can approach the Employment Compensation Commissioner to seek compensation.
In March 2016, Hindustan Unilever entered into an out of court settlement with its ex-employees to provide "undisclosed" ex-gratia payment, in addition to long-term health and well-being benefits, to 511 of its former workers of the thermometer factory who were exposed to toxic mercury vapour. Accordingly, the ex-employees withdrew the 'class action litigation' before the Madras High Court and the High Court of Justice, London.
References
Mercury poisoning
Bioremediation
Ecological restoration
Kodaikanal
Neurotoxins
Pollution in India
Environment of Tamil Nadu
Environmental justice
Environmental disasters in India
Minamata disease
Environmental history of India | Kodaikanal mercury poisoning | [
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40,409,788 | https://en.wikipedia.org/wiki/Convolutional%20neural%20network | A convolutional neural network (CNN) is a regularized type of feed-forward neural network that learns features by itself via filter (or kernel) optimization. This type of deep learning network has been applied to process and make predictions from many different types of data including text, images and audio. Convolution-based networks are the de-facto standard in deep learning-based approaches to computer vision and image processing, and have only recently been replaced—in some cases—by newer deep learning architectures such as the transformer. Vanishing gradients and exploding gradients, seen during backpropagation in earlier neural networks, are prevented by using regularized weights over fewer connections. For example, for each neuron in the fully-connected layer, 10,000 weights would be required for processing an image sized 100 × 100 pixels. However, applying cascaded convolution (or cross-correlation) kernels, only 25 neurons are required to process 5x5-sized tiles. Higher-layer features are extracted from wider context windows, compared to lower-layer features.
Some applications of CNNs include:
image and video recognition,
recommender systems,
image classification,
image segmentation,
medical image analysis,
natural language processing,
brain–computer interfaces, and
financial time series.
CNNs are also known as shift invariant or space invariant artificial neural networks, based on the shared-weight architecture of the convolution kernels or filters that slide along input features and provide translation-equivariant responses known as feature maps. Counter-intuitively, most convolutional neural networks are not invariant to translation, due to the downsampling operation they apply to the input.
Feed-forward neural networks are usually fully connected networks, that is, each neuron in one layer is connected to all neurons in the next layer. The "full connectivity" of these networks makes them prone to overfitting data. Typical ways of regularization, or preventing overfitting, include: penalizing parameters during training (such as weight decay) or trimming connectivity (skipped connections, dropout, etc.) Robust datasets also increase the probability that CNNs will learn the generalized principles that characterize a given dataset rather than the biases of a poorly-populated set.
Convolutional networks were inspired by biological processes in that the connectivity pattern between neurons resembles the organization of the animal visual cortex. Individual cortical neurons respond to stimuli only in a restricted region of the visual field known as the receptive field. The receptive fields of different neurons partially overlap such that they cover the entire visual field.
CNNs use relatively little pre-processing compared to other image classification algorithms. This means that the network learns to optimize the filters (or kernels) through automated learning, whereas in traditional algorithms these filters are hand-engineered. This independence from prior knowledge and human intervention in feature extraction is a major advantage.
Architecture
A convolutional neural network consists of an input layer, hidden layers and an output layer. In a convolutional neural network, the hidden layers include one or more layers that perform convolutions. Typically this includes a layer that performs a dot product of the convolution kernel with the layer's input matrix. This product is usually the Frobenius inner product, and its activation function is commonly ReLU. As the convolution kernel slides along the input matrix for the layer, the convolution operation generates a feature map, which in turn contributes to the input of the next layer. This is followed by other layers such as pooling layers, fully connected layers, and normalization layers.
Here it should be noted how close a convolutional neural network is to a matched filter.
Convolutional layers
In a CNN, the input is a tensor with shape:
(number of inputs) × (input height) × (input width) × (input channels)
After passing through a convolutional layer, the image becomes abstracted to a feature map, also called an activation map, with shape:
(number of inputs) × (feature map height) × (feature map width) × (feature map channels).
Convolutional layers convolve the input and pass its result to the next layer. This is similar to the response of a neuron in the visual cortex to a specific stimulus. Each convolutional neuron processes data only for its receptive field.
Although fully connected feedforward neural networks can be used to learn features and classify data, this architecture is generally impractical for larger inputs (e.g., high-resolution images), which would require massive numbers of neurons because each pixel is a relevant input feature. A fully connected layer for an image of size 100 × 100 has 10,000 weights for each neuron in the second layer. Convolution reduces the number of free parameters, allowing the network to be deeper. For example, using a 5 × 5 tiling region, each with the same shared weights, requires only 25 neurons. Using regularized weights over fewer parameters avoids the vanishing gradients and exploding gradients problems seen during backpropagation in earlier neural networks.
To speed processing, standard convolutional layers can be replaced by depthwise separable convolutional layers, which are based on a depthwise convolution followed by a pointwise convolution. The depthwise convolution is a spatial convolution applied independently over each channel of the input tensor, while the pointwise convolution is a standard convolution restricted to the use of kernels.
Pooling layers
Convolutional networks may include local and/or global pooling layers along with traditional convolutional layers. Pooling layers reduce the dimensions of data by combining the outputs of neuron clusters at one layer into a single neuron in the next layer. Local pooling combines small clusters, tiling sizes such as 2 × 2 are commonly used. Global pooling acts on all the neurons of the feature map. There are two common types of pooling in popular use: max and average. Max pooling uses the maximum value of each local cluster of neurons in the feature map, while average pooling takes the average value.
Fully connected layers
Fully connected layers connect every neuron in one layer to every neuron in another layer. It is the same as a traditional multilayer perceptron neural network (MLP). The flattened matrix goes through a fully connected layer to classify the images.
Receptive field
In neural networks, each neuron receives input from some number of locations in the previous layer. In a convolutional layer, each neuron receives input from only a restricted area of the previous layer called the neuron's receptive field. Typically the area is a square (e.g. 5 by 5 neurons). Whereas, in a fully connected layer, the receptive field is the entire previous layer. Thus, in each convolutional layer, each neuron takes input from a larger area in the input than previous layers. This is due to applying the convolution over and over, which takes the value of a pixel into account, as well as its surrounding pixels. When using dilated layers, the number of pixels in the receptive field remains constant, but the field is more sparsely populated as its dimensions grow when combining the effect of several layers.
To manipulate the receptive field size as desired, there are some alternatives to the standard convolutional layer. For example, atrous or dilated convolution expands the receptive field size without increasing the number of parameters by interleaving visible and blind regions. Moreover, a single dilated convolutional layer can comprise filters with multiple dilation ratios, thus having a variable receptive field size.
Weights
Each neuron in a neural network computes an output value by applying a specific function to the input values received from the receptive field in the previous layer. The function that is applied to the input values is determined by a vector of weights and a bias (typically real numbers). Learning consists of iteratively adjusting these biases and weights.
The vectors of weights and biases are called filters and represent particular features of the input (e.g., a particular shape). A distinguishing feature of CNNs is that many neurons can share the same filter. This reduces the memory footprint because a single bias and a single vector of weights are used across all receptive fields that share that filter, as opposed to each receptive field having its own bias and vector weighting.
Deconvolutional
A deconvolutional neural network is essentially the reverse of a CNN. It consists of deconvolutional layers and unpooling layers.
A deconvolutional layer is the transpose of a convolutional layer. Specifically, a convolutional layer can be written as a multiplication with a matrix, and a deconvolutional layer is multiplication with the transpose of that matrix.
An unpooling layer expands the layer. The max-unpooling layer is the simplest, as it simply copies each entry multiple times. For example, a 2-by-2 max-unpooling layer is .
Deconvolution layers are used in image generators. By default, it creates periodic checkerboard artifact, which can be fixed by upscale-then-convolve.
History
CNN are often compared to the way the brain achieves vision processing in living organisms.
Receptive fields in the visual cortex
Work by Hubel and Wiesel in the 1950s and 1960s showed that cat visual cortices contain neurons that individually respond to small regions of the visual field. Provided the eyes are not moving, the region of visual space within which visual stimuli affect the firing of a single neuron is known as its receptive field. Neighboring cells have similar and overlapping receptive fields. Receptive field size and location varies systematically across the cortex to form a complete map of visual space. The cortex in each hemisphere represents the contralateral visual field.
Their 1968 paper identified two basic visual cell types in the brain:
simple cells, whose output is maximized by straight edges having particular orientations within their receptive field
complex cells, which have larger receptive fields, whose output is insensitive to the exact position of the edges in the field.
Hubel and Wiesel also proposed a cascading model of these two types of cells for use in pattern recognition tasks.
Neocognitron, origin of the CNN architecture
Inspired by Hubel and Wiesel's work, in 1969, Kunihiko Fukushima published a deep CNN that uses ReLU activation function. Unlike most modern networks, this network used hand-designed kernels. It was not used in his neocognitron, since all the weights were nonnegative; lateral inhibition was used instead. The rectifier has become the most popular activation function for CNNs and deep neural networks in general.
The "neocognitron" was introduced by Kunihiko Fukushima in 1979. The kernels were trained by unsupervised learning. It was inspired by the above-mentioned work of Hubel and Wiesel. The neocognitron introduced the two basic types of layers:
"S-layer": a shared-weights receptive-field layer, later known as a convolutional layer, which contains units whose receptive fields cover a patch of the previous layer. A shared-weights receptive-field group (a "plane" in neocognitron terminology) is often called a filter, and a layer typically has several such filters.
"C-layer": a downsampling layer that contain units whose receptive fields cover patches of previous convolutional layers. Such a unit typically computes a weighted average of the activations of the units in its patch, and applies inhibition (divisive normalization) pooled from a somewhat larger patch and across different filters in a layer, and applies a saturating activation function. The patch weights are nonnegative and are not trainable in the original neocognitron. The downsampling and competitive inhibition help to classify features and objects in visual scenes even when the objects are shifted.
In a variant of the neocognitron called the cresceptron, instead of using Fukushima's spatial averaging with inhibition and saturation, J. Weng et al. in 1993 introduced a method called max-pooling where a downsampling unit computes the maximum of the activations of the units in its patch. Max-pooling is often used in modern CNNs.
Several supervised and unsupervised learning algorithms have been proposed over the decades to train the weights of a neocognitron. Today, however, the CNN architecture is usually trained through backpropagation.
Convolution in time
The term "convolution" first appears in neural networks in a paper by Toshiteru Homma, Les Atlas, and Robert Marks II at the first Conference on Neural Information Processing Systems in 1987. Their paper replaced multiplication with convolution in time, inherently providing shift invariance, motivated by and connecting more directly to the signal-processing concept of a filter, and demonstrated it on a speech recognition task. They also pointed out that as a data-trainable system, convolution is essentially equivalent to correlation since reversal of the weights does not affect the final learned function ("For convenience, we denote * as correlation instead of convolution. Note that convolving a(t) with b(t) is equivalent to correlating a(-t) with b(t)."). Modern CNN implementations typically do correlation and call it convolution, for convenience, as they did here.
Time delay neural networks
The time delay neural network (TDNN) was introduced in 1987 by Alex Waibel et al. for phoneme recognition and was one of the first convolutional networks, as it achieved shift-invariance. A TDNN is a 1-D convolutional neural net where the convolution is performed along the time axis of the data. It is the first CNN utilizing weight sharing in combination with a training by gradient descent, using backpropagation. Thus, while also using a pyramidal structure as in the neocognitron, it performed a global optimization of the weights instead of a local one.
TDNNs are convolutional networks that share weights along the temporal dimension. They allow speech signals to be processed time-invariantly. In 1990 Hampshire and Waibel introduced a variant that performs a two-dimensional convolution. Since these TDNNs operated on spectrograms, the resulting phoneme recognition system was invariant to both time and frequency shifts, as with images processed by a neocognitron.
TDNNs improved the performance of far-distance speech recognition.
Image recognition with CNNs trained by gradient descent
Denker et al. (1989) designed a 2-D CNN system to recognize hand-written ZIP Code numbers. However, the lack of an efficient training method to determine the kernel coefficients of the involved convolutions meant that all the coefficients had to be laboriously hand-designed.
Following the advances in the training of 1-D CNNs by Waibel et al. (1987), Yann LeCun et al. (1989) used back-propagation to learn the convolution kernel coefficients directly from images of hand-written numbers. Learning was thus fully automatic, performed better than manual coefficient design, and was suited to a broader range of image recognition problems and image types.
Wei Zhang et al. (1988) used back-propagation to train the convolution kernels of a CNN for alphabets recognition. The model was called shift-invariant pattern recognition neural network before the name CNN was coined later in the early 1990s. Wei Zhang et al. also applied the same CNN without the last fully connected layer for medical image object segmentation (1991) and breast cancer detection in mammograms (1994).
This approach became a foundation of modern computer vision.
Max pooling
In 1990 Yamaguchi et al. introduced the concept of max pooling, a fixed filtering operation that calculates and propagates the maximum value of a given region. They did so by combining TDNNs with max pooling to realize a speaker-independent isolated word recognition system. In their system they used several TDNNs per word, one for each syllable. The results of each TDNN over the input signal were combined using max pooling and the outputs of the pooling layers were then passed on to networks performing the actual word classification.
LeNet-5
LeNet-5, a pioneering 7-level convolutional network by LeCun et al. in 1995, classifies hand-written numbers on checks () digitized in 32x32 pixel images. The ability to process higher-resolution images requires larger and more layers of convolutional neural networks, so this technique is constrained by the availability of computing resources.
It was superior than other commercial courtesy amount reading systems (as of 1995). The system was integrated in NCR's check reading systems, and fielded in several American banks since June 1996, reading millions of checks per day.
Shift-invariant neural network
A shift-invariant neural network was proposed by Wei Zhang et al. for image character recognition in 1988. It is a modified Neocognitron by keeping only the convolutional interconnections between the image feature layers and the last fully connected layer. The model was trained with back-propagation. The training algorithm was further improved in 1991 to improve its generalization ability. The model architecture was modified by removing the last fully connected layer and applied for medical image segmentation (1991) and automatic detection of breast cancer in mammograms (1994).
A different convolution-based design was proposed in 1988 for application to decomposition of one-dimensional electromyography convolved signals via de-convolution. This design was modified in 1989 to other de-convolution-based designs.
Topological deep learning
Topological deep learning was first introduced in 2017.
It integrates topological data analysis and convolutional neural networks for intricately complex data.
Topological deep learning has become a new frontier in deep learning.
GPU implementations
Although CNNs were invented in the 1980s, their breakthrough in the 2000s required fast implementations on graphics processing units (GPUs).
In 2004, it was shown by K. S. Oh and K. Jung that standard neural networks can be greatly accelerated on GPUs. Their implementation was 20 times faster than an equivalent implementation on CPU. In 2005, another paper also emphasised the value of GPGPU for machine learning.
The first GPU-implementation of a CNN was described in 2006 by K. Chellapilla et al. Their implementation was 4 times faster than an equivalent implementation on CPU. In the same period, GPUs were also used for unsupervised training of deep belief networks.
In 2010, Dan Ciresan et al. at IDSIA trained deep feedforward networks on GPUs. In 2011, they extended this to CNNs, accelerating by 60 compared to training CPU. In 2011, the network win an image recognition contest where they achieved superhuman performance for the first time. Then they won more competitions and achieved state of the art on several benchmarks.
Subsequently, AlexNet, a similar GPU-based CNN by Alex Krizhevsky et al. won the ImageNet Large Scale Visual Recognition Challenge 2012. It was an early catalytic event for the AI boom.
Compared to the training of CNNs using GPUs, not much attention was given to CPU. (Viebke et al 2019) parallelizes CNN by thread- and SIMD-level parallelism that is available on the Intel Xeon Phi.
Distinguishing features
In the past, traditional multilayer perceptron (MLP) models were used for image recognition. However, the full connectivity between nodes caused the curse of dimensionality, and was computationally intractable with higher-resolution images. A 1000×1000-pixel image with RGB color channels has 3 million weights per fully-connected neuron, which is too high to feasibly process efficiently at scale.
For example, in CIFAR-10, images are only of size 32×32×3 (32 wide, 32 high, 3 color channels), so a single fully connected neuron in the first hidden layer of a regular neural network would have 32*32*3 = 3,072 weights. A 200×200 image, however, would lead to neurons that have 200*200*3 = 120,000 weights.
Also, such network architecture does not take into account the spatial structure of data, treating input pixels which are far apart in the same way as pixels that are close together. This ignores locality of reference in data with a grid-topology (such as images), both computationally and semantically. Thus, full connectivity of neurons is wasteful for purposes such as image recognition that are dominated by spatially local input patterns.
Convolutional neural networks are variants of multilayer perceptrons, designed to emulate the behavior of a visual cortex. These models mitigate the challenges posed by the MLP architecture by exploiting the strong spatially local correlation present in natural images. As opposed to MLPs, CNNs have the following distinguishing features:
3D volumes of neurons. The layers of a CNN have neurons arranged in 3 dimensions: width, height and depth. Where each neuron inside a convolutional layer is connected to only a small region of the layer before it, called a receptive field. Distinct types of layers, both locally and completely connected, are stacked to form a CNN architecture.
Local connectivity: following the concept of receptive fields, CNNs exploit spatial locality by enforcing a local connectivity pattern between neurons of adjacent layers. The architecture thus ensures that the learned "filters" produce the strongest response to a spatially local input pattern. Stacking many such layers leads to nonlinear filters that become increasingly global (i.e. responsive to a larger region of pixel space) so that the network first creates representations of small parts of the input, then from them assembles representations of larger areas.
Shared weights: In CNNs, each filter is replicated across the entire visual field. These replicated units share the same parameterization (weight vector and bias) and form a feature map. This means that all the neurons in a given convolutional layer respond to the same feature within their specific response field. Replicating units in this way allows for the resulting activation map to be equivariant under shifts of the locations of input features in the visual field, i.e. they grant translational equivariance—given that the layer has a stride of one.
Pooling: In a CNN's pooling layers, feature maps are divided into rectangular sub-regions, and the features in each rectangle are independently down-sampled to a single value, commonly by taking their average or maximum value. In addition to reducing the sizes of feature maps, the pooling operation grants a degree of local translational invariance to the features contained therein, allowing the CNN to be more robust to variations in their positions.
Together, these properties allow CNNs to achieve better generalization on vision problems. Weight sharing dramatically reduces the number of free parameters learned, thus lowering the memory requirements for running the network and allowing the training of larger, more powerful networks.
Building blocks
A CNN architecture is formed by a stack of distinct layers that transform the input volume into an output volume (e.g. holding the class scores) through a differentiable function. A few distinct types of layers are commonly used. These are further discussed below.
Convolutional layer
The convolutional layer is the core building block of a CNN. The layer's parameters consist of a set of learnable filters (or kernels), which have a small receptive field, but extend through the full depth of the input volume. During the forward pass, each filter is convolved across the width and height of the input volume, computing the dot product between the filter entries and the input, producing a 2-dimensional activation map of that filter. As a result, the network learns filters that activate when it detects some specific type of feature at some spatial position in the input.
Stacking the activation maps for all filters along the depth dimension forms the full output volume of the convolution layer. Every entry in the output volume can thus also be interpreted as an output of a neuron that looks at a small region in the input. Each entry in an activation map use the same set of parameters that define the filter.
Self-supervised learning has been adapted for use in convolutional layers by using sparse patches with a high-mask ratio and a global response normalization layer.
Local connectivity
When dealing with high-dimensional inputs such as images, it is impractical to connect neurons to all neurons in the previous volume because such a network architecture does not take the spatial structure of the data into account. Convolutional networks exploit spatially local correlation by enforcing a sparse local connectivity pattern between neurons of adjacent layers: each neuron is connected to only a small region of the input volume.
The extent of this connectivity is a hyperparameter called the receptive field of the neuron. The connections are local in space (along width and height), but always extend along the entire depth of the input volume. Such an architecture ensures that the learned () filters produce the strongest response to a spatially local input pattern.
Spatial arrangement
Three hyperparameters control the size of the output volume of the convolutional layer: the depth, stride, and padding size:
The depth of the output volume controls the number of neurons in a layer that connect to the same region of the input volume. These neurons learn to activate for different features in the input. For example, if the first convolutional layer takes the raw image as input, then different neurons along the depth dimension may activate in the presence of various oriented edges, or blobs of color.
Stride controls how depth columns around the width and height are allocated. If the stride is 1, then we move the filters one pixel at a time. This leads to heavily overlapping receptive fields between the columns, and to large output volumes. For any integer a stride S means that the filter is translated S units at a time per output. In practice, is rare. A greater stride means smaller overlap of receptive fields and smaller spatial dimensions of the output volume.
Sometimes, it is convenient to pad the input with zeros (or other values, such as the average of the region) on the border of the input volume. The size of this padding is a third hyperparameter. Padding provides control of the output volume's spatial size. In particular, sometimes it is desirable to exactly preserve the spatial size of the input volume, this is commonly referred to as "same" padding.
The spatial size of the output volume is a function of the input volume size , the kernel field size of the convolutional layer neurons, the stride , and the amount of zero padding on the border. The number of neurons that "fit" in a given volume is then:
If this number is not an integer, then the strides are incorrect and the neurons cannot be tiled to fit across the input volume in a symmetric way. In general, setting zero padding to be when the stride is ensures that the input volume and output volume will have the same size spatially. However, it is not always completely necessary to use all of the neurons of the previous layer. For example, a neural network designer may decide to use just a portion of padding.
Parameter sharing
A parameter sharing scheme is used in convolutional layers to control the number of free parameters. It relies on the assumption that if a patch feature is useful to compute at some spatial position, then it should also be useful to compute at other positions. Denoting a single 2-dimensional slice of depth as a depth slice, the neurons in each depth slice are constrained to use the same weights and bias.
Since all neurons in a single depth slice share the same parameters, the forward pass in each depth slice of the convolutional layer can be computed as a convolution of the neuron's weights with the input volume. Therefore, it is common to refer to the sets of weights as a filter (or a kernel), which is convolved with the input. The result of this convolution is an activation map, and the set of activation maps for each different filter are stacked together along the depth dimension to produce the output volume. Parameter sharing contributes to the translation invariance of the CNN architecture.
Sometimes, the parameter sharing assumption may not make sense. This is especially the case when the input images to a CNN have some specific centered structure; for which we expect completely different features to be learned on different spatial locations. One practical example is when the inputs are faces that have been centered in the image: we might expect different eye-specific or hair-specific features to be learned in different parts of the image. In that case it is common to relax the parameter sharing scheme, and instead simply call the layer a "locally connected layer".
Pooling layer
Another important concept of CNNs is pooling, which is used as a form of non-linear down-sampling. Pooling provides downsampling because it reduces the spatial dimensions (height and width) of the input feature maps while retaining the most important information. There are several non-linear functions to implement pooling, where max pooling and average pooling are the most common. Pooling aggregates information from small regions of the input creating partitions of the input feature map, typically using a fixed-size window (like 2x2) and applying a stride (often 2) to move the window across the input. Note that without using a stride greater than 1, pooling would not perform downsampling, as it would simply move the pooling window across the input one step at a time, without reducing the size of the feature map. In other words, the stride is what actually causes the downsampling by determining how much the pooling window moves over the input.
Intuitively, the exact location of a feature is less important than its rough location relative to other features. This is the idea behind the use of pooling in convolutional neural networks. The pooling layer serves to progressively reduce the spatial size of the representation, to reduce the number of parameters, memory footprint and amount of computation in the network, and hence to also control overfitting. This is known as down-sampling. It is common to periodically insert a pooling layer between successive convolutional layers (each one typically followed by an activation function, such as a ReLU layer) in a CNN architecture. While pooling layers contribute to local translation invariance, they do not provide global translation invariance in a CNN, unless a form of global pooling is used. The pooling layer commonly operates independently on every depth, or slice, of the input and resizes it spatially. A very common form of max pooling is a layer with filters of size 2×2, applied with a stride of 2, which subsamples every depth slice in the input by 2 along both width and height, discarding 75% of the activations:
In this case, every max operation is over 4 numbers. The depth dimension remains unchanged (this is true for other forms of pooling as well).
In addition to max pooling, pooling units can use other functions, such as average pooling or ℓ2-norm pooling. Average pooling was often used historically but has recently fallen out of favor compared to max pooling, which generally performs better in practice.
Due to the effects of fast spatial reduction of the size of the representation, there is a recent trend towards using smaller filters or discarding pooling layers altogether.
Channel max pooling
A channel max pooling (CMP) operation layer conducts the MP operation along the channel side among the corresponding positions of the consecutive feature maps for the purpose of redundant information elimination. The CMP makes the significant features gather together within fewer channels, which is important for fine-grained image classification that needs more discriminating features. Meanwhile, another advantage of the CMP operation is to make the channel number of feature maps smaller before it connects to the first fully connected (FC) layer. Similar to the MP operation, we denote the input feature maps and output feature maps of a CMP layer as F ∈ R(C×M×N) and C ∈ R(c×M×N), respectively, where C and c are the channel numbers of the input and output feature maps, M and N are the widths and the height of the feature maps, respectively. Note that the CMP operation only changes the channel number of the feature maps. The width and the height of the feature maps are not changed, which is different from the MP operation.
See for reviews for pooling methods.
ReLU layer
ReLU is the abbreviation of rectified linear unit. It was proposed by Alston Householder in 1941, and used in CNN by Kunihiko Fukushima in 1969. ReLU applies the non-saturating activation function . It effectively removes negative values from an activation map by setting them to zero. It introduces nonlinearity to the decision function and in the overall network without affecting the receptive fields of the convolution layers.
In 2011, Xavier Glorot, Antoine Bordes and Yoshua Bengio found that ReLU enables better training of deeper networks, compared to widely used activation functions prior to 2011.
Other functions can also be used to increase nonlinearity, for example the saturating hyperbolic tangent , , and the sigmoid function . ReLU is often preferred to other functions because it trains the neural network several times faster without a significant penalty to generalization accuracy.
Fully connected layer
After several convolutional and max pooling layers, the final classification is done via fully connected layers. Neurons in a fully connected layer have connections to all activations in the previous layer, as seen in regular (non-convolutional) artificial neural networks. Their activations can thus be computed as an affine transformation, with matrix multiplication followed by a bias offset (vector addition of a learned or fixed bias term).
Loss layer
The "loss layer", or "loss function", specifies how training penalizes the deviation between the predicted output of the network, and the true data labels (during supervised learning). Various loss functions can be used, depending on the specific task.
The Softmax loss function is used for predicting a single class of K mutually exclusive classes. Sigmoid cross-entropy loss is used for predicting K independent probability values in . Euclidean loss is used for regressing to real-valued labels .
Hyperparameters
Hyperparameters are various settings that are used to control the learning process. CNNs use more hyperparameters than a standard multilayer perceptron (MLP).
Kernel size
The kernel is the number of pixels processed together. It is typically expressed as the kernel's dimensions, e.g., 2x2, or 3x3.
Padding
Padding is the addition of (typically) 0-valued pixels on the borders of an image. This is done so that the border pixels are not undervalued (lost) from the output because they would ordinarily participate in only a single receptive field instance. The padding applied is typically one less than the corresponding kernel dimension. For example, a convolutional layer using 3x3 kernels would receive a 2-pixel pad, that is 1 pixel on each side of the image.
Stride
The stride is the number of pixels that the analysis window moves on each iteration. A stride of 2 means that each kernel is offset by 2 pixels from its predecessor.
Number of filters
Since feature map size decreases with depth, layers near the input layer tend to have fewer filters while higher layers can have more. To equalize computation at each layer, the product of feature values va with pixel position is kept roughly constant across layers. Preserving more information about the input would require keeping the total number of activations (number of feature maps times number of pixel positions) non-decreasing from one layer to the next.
The number of feature maps directly controls the capacity and depends on the number of available examples and task complexity.
Filter size
Common filter sizes found in the literature vary greatly, and are usually chosen based on the data set. Typical filter sizes range from 1x1 to 7x7. As two famous examples, AlexNet used 3x3, 5x5, and 11x11. Inceptionv3 used 1x1, 3x3, and 5x5.
The challenge is to find the right level of granularity so as to create abstractions at the proper scale, given a particular data set, and without overfitting.
Pooling type and size
Max pooling is typically used, often with a 2x2 dimension. This implies that the input is drastically downsampled, reducing processing cost.
Greater pooling reduces the dimension of the signal, and may result in unacceptable information loss. Often, non-overlapping pooling windows perform best.
Dilation
Dilation involves ignoring pixels within a kernel. This reduces processing memory potentially without significant signal loss. A dilation of 2 on a 3x3 kernel expands the kernel to 5x5, while still processing 9 (evenly spaced) pixels. Specifically, the processed pixels after the dilation are the cells (1,1), (1,3), (1,5), (3,1), (3,3), (3,5), (5,1), (5,3), (5,5), where (i,j) denotes the cell of the i-th row and j-th column in the expanded 5x5 kernel. Accordingly, dilation of 4 expands the kernel to 7x7.
Translation equivariance and aliasing
It is commonly assumed that CNNs are invariant to shifts of the input. Convolution or pooling layers within a CNN that do not have a stride greater than one are indeed equivariant to translations of the input. However, layers with a stride greater than one ignore the Nyquist-Shannon sampling theorem and might lead to aliasing of the input signal While, in principle, CNNs are capable of implementing anti-aliasing filters, it has been observed that this does not happen in practice and yield models that are not equivariant to translations.
Furthermore, if a CNN makes use of fully connected layers, translation equivariance does not imply translation invariance, as the fully connected layers are not invariant to shifts of the input. One solution for complete translation invariance is avoiding any down-sampling throughout the network and applying global average pooling at the last layer. Additionally, several other partial solutions have been proposed, such as anti-aliasing before downsampling operations, spatial transformer networks, data augmentation, subsampling combined with pooling, and capsule neural networks.
Evaluation
The accuracy of the final model is based on a sub-part of the dataset set apart at the start, often called a test-set. Other times methods such as k-fold cross-validation are applied. Other strategies include using conformal prediction.
Regularization methods
Regularization is a process of introducing additional information to solve an ill-posed problem or to prevent overfitting. CNNs use various types of regularization.
Empirical
Dropout
Because a fully connected layer occupies most of the parameters, it is prone to overfitting. One method to reduce overfitting is dropout, introduced in 2014. At each training stage, individual nodes are either "dropped out" of the net (ignored) with probability or kept with probability , so that a reduced network is left; incoming and outgoing edges to a dropped-out node are also removed. Only the reduced network is trained on the data in that stage. The removed nodes are then reinserted into the network with their original weights.
In the training stages, is usually 0.5; for input nodes, it is typically much higher because information is directly lost when input nodes are ignored.
At testing time after training has finished, we would ideally like to find a sample average of all possible dropped-out networks; unfortunately this is unfeasible for large values of . However, we can find an approximation by using the full network with each node's output weighted by a factor of , so the expected value of the output of any node is the same as in the training stages. This is the biggest contribution of the dropout method: although it effectively generates neural nets, and as such allows for model combination, at test time only a single network needs to be tested.
By avoiding training all nodes on all training data, dropout decreases overfitting. The method also significantly improves training speed. This makes the model combination practical, even for deep neural networks. The technique seems to reduce node interactions, leading them to learn more robust features that better generalize to new data.
DropConnect
DropConnect is the generalization of dropout in which each connection, rather than each output unit, can be dropped with probability . Each unit thus receives input from a random subset of units in the previous layer.
DropConnect is similar to dropout as it introduces dynamic sparsity within the model, but differs in that the sparsity is on the weights, rather than the output vectors of a layer. In other words, the fully connected layer with DropConnect becomes a sparsely connected layer in which the connections are chosen at random during the training stage.
Stochastic pooling
A major drawback to Dropout is that it does not have the same benefits for convolutional layers, where the neurons are not fully connected.
Even before Dropout, in 2013 a technique called stochastic pooling, the conventional deterministic pooling operations were replaced with a stochastic procedure, where the activation within each pooling region is picked randomly according to a multinomial distribution, given by the activities within the pooling region. This approach is free of hyperparameters and can be combined with other regularization approaches, such as dropout and data augmentation.
An alternate view of stochastic pooling is that it is equivalent to standard max pooling but with many copies of an input image, each having small local deformations. This is similar to explicit elastic deformations of the input images, which delivers excellent performance on the MNIST data set. Using stochastic pooling in a multilayer model gives an exponential number of deformations since the selections in higher layers are independent of those below.
Artificial data
Because the degree of model overfitting is determined by both its power and the amount of training it receives, providing a convolutional network with more training examples can reduce overfitting. Because there is often not enough available data to train, especially considering that some part should be spared for later testing, two approaches are to either generate new data from scratch (if possible) or perturb existing data to create new ones. The latter one is used since mid-1990s. For example, input images can be cropped, rotated, or rescaled to create new examples with the same labels as the original training set.
Explicit
Early stopping
One of the simplest methods to prevent overfitting of a network is to simply stop the training before overfitting has had a chance to occur. It comes with the disadvantage that the learning process is halted.
Number of parameters
Another simple way to prevent overfitting is to limit the number of parameters, typically by limiting the number of hidden units in each layer or limiting network depth. For convolutional networks, the filter size also affects the number of parameters. Limiting the number of parameters restricts the predictive power of the network directly, reducing the complexity of the function that it can perform on the data, and thus limits the amount of overfitting. This is equivalent to a "zero norm".
Weight decay
A simple form of added regularizer is weight decay, which simply adds an additional error, proportional to the sum of weights (L1 norm) or squared magnitude (L2 norm) of the weight vector, to the error at each node. The level of acceptable model complexity can be reduced by increasing the proportionality constant('alpha' hyperparameter), thus increasing the penalty for large weight vectors.
L2 regularization is the most common form of regularization. It can be implemented by penalizing the squared magnitude of all parameters directly in the objective. The L2 regularization has the intuitive interpretation of heavily penalizing peaky weight vectors and preferring diffuse weight vectors. Due to multiplicative interactions between weights and inputs this has the useful property of encouraging the network to use all of its inputs a little rather than some of its inputs a lot.
L1 regularization is also common. It makes the weight vectors sparse during optimization. In other words, neurons with L1 regularization end up using only a sparse subset of their most important inputs and become nearly invariant to the noisy inputs. L1 with L2 regularization can be combined; this is called elastic net regularization.
Max norm constraints
Another form of regularization is to enforce an absolute upper bound on the magnitude of the weight vector for every neuron and use projected gradient descent to enforce the constraint. In practice, this corresponds to performing the parameter update as normal, and then enforcing the constraint by clamping the weight vector of every neuron to satisfy . Typical values of are order of 3–4. Some papers report improvements when using this form of regularization.
Hierarchical coordinate frames
Pooling loses the precise spatial relationships between high-level parts (such as nose and mouth in a face image). These relationships are needed for identity recognition. Overlapping the pools so that each feature occurs in multiple pools, helps retain the information. Translation alone cannot extrapolate the understanding of geometric relationships to a radically new viewpoint, such as a different orientation or scale. On the other hand, people are very good at extrapolating; after seeing a new shape once they can recognize it from a different viewpoint.
An earlier common way to deal with this problem is to train the network on transformed data in different orientations, scales, lighting, etc. so that the network can cope with these variations. This is computationally intensive for large data-sets. The alternative is to use a hierarchy of coordinate frames and use a group of neurons to represent a conjunction of the shape of the feature and its pose relative to the retina. The pose relative to the retina is the relationship between the coordinate frame of the retina and the intrinsic features' coordinate frame.
Thus, one way to represent something is to embed the coordinate frame within it. This allows large features to be recognized by using the consistency of the poses of their parts (e.g. nose and mouth poses make a consistent prediction of the pose of the whole face). This approach ensures that the higher-level entity (e.g. face) is present when the lower-level (e.g. nose and mouth) agree on its prediction of the pose. The vectors of neuronal activity that represent pose ("pose vectors") allow spatial transformations modeled as linear operations that make it easier for the network to learn the hierarchy of visual entities and generalize across viewpoints. This is similar to the way the human visual system imposes coordinate frames in order to represent shapes.
Applications
Image recognition
CNNs are often used in image recognition systems. In 2012, an error rate of 0.23% on the MNIST database was reported. Another paper on using CNN for image classification reported that the learning process was "surprisingly fast"; in the same paper, the best published results as of 2011 were achieved in the MNIST database and the NORB database. Subsequently, a similar CNN called AlexNet won the ImageNet Large Scale Visual Recognition Challenge 2012.
When applied to facial recognition, CNNs achieved a large decrease in error rate. Another paper reported a 97.6% recognition rate on "5,600 still images of more than 10 subjects". CNNs were used to assess video quality in an objective way after manual training; the resulting system had a very low root mean square error.
The ImageNet Large Scale Visual Recognition Challenge is a benchmark in object classification and detection, with millions of images and hundreds of object classes. In the ILSVRC 2014, a large-scale visual recognition challenge, almost every highly ranked team used CNN as their basic framework. The winner GoogLeNet (the foundation of DeepDream) increased the mean average precision of object detection to 0.439329, and reduced classification error to 0.06656, the best result to date. Its network applied more than 30 layers. That performance of convolutional neural networks on the ImageNet tests was close to that of humans. The best algorithms still struggle with objects that are small or thin, such as a small ant on a stem of a flower or a person holding a quill in their hand. They also have trouble with images that have been distorted with filters, an increasingly common phenomenon with modern digital cameras. By contrast, those kinds of images rarely trouble humans. Humans, however, tend to have trouble with other issues. For example, they are not good at classifying objects into fine-grained categories such as the particular breed of dog or species of bird, whereas convolutional neural networks handle this.
In 2015, a many-layered CNN demonstrated the ability to spot faces from a wide range of angles, including upside down, even when partially occluded, with competitive performance. The network was trained on a database of 200,000 images that included faces at various angles and orientations and a further 20 million images without faces. They used batches of 128 images over 50,000 iterations.
Video analysis
Compared to image data domains, there is relatively little work on applying CNNs to video classification. Video is more complex than images since it has another (temporal) dimension. However, some extensions of CNNs into the video domain have been explored. One approach is to treat space and time as equivalent dimensions of the input and perform convolutions in both time and space. Another way is to fuse the features of two convolutional neural networks, one for the spatial and one for the temporal stream. Long short-term memory (LSTM) recurrent units are typically incorporated after the CNN to account for inter-frame or inter-clip dependencies. Unsupervised learning schemes for training spatio-temporal features have been introduced, based on Convolutional Gated Restricted Boltzmann Machines and Independent Subspace Analysis. Its application can be seen in text-to-video model.
Natural language processing
CNNs have also been explored for natural language processing. CNN models are effective for various NLP problems and achieved excellent results in semantic parsing, search query retrieval, sentence modeling, classification, prediction and other traditional NLP tasks.
Compared to traditional language processing methods such as recurrent neural networks, CNNs can represent different contextual realities of language that do not rely on a series-sequence assumption, while RNNs are better suitable when classical time series modeling is required.
Anomaly detection
A CNN with 1-D convolutions was used on time series in the frequency domain (spectral residual) by an unsupervised model to detect anomalies in the time domain.
Drug discovery
CNNs have been used in drug discovery. Predicting the interaction between molecules and biological proteins can identify potential treatments. In 2015, Atomwise introduced AtomNet, the first deep learning neural network for structure-based drug design. The system trains directly on 3-dimensional representations of chemical interactions. Similar to how image recognition networks learn to compose smaller, spatially proximate features into larger, complex structures, AtomNet discovers chemical features, such as aromaticity, sp3 carbons, and hydrogen bonding. Subsequently, AtomNet was used to predict novel candidate biomolecules for multiple disease targets, most notably treatments for the Ebola virus and multiple sclerosis.
In 2016-2019, topologyNet and mathematical deep learning
achieved first place in multiple categories of the D3R Grand Challenges, a worldwide annual competition series focused on computer-aided drug design.
Checkers game
CNNs have been used in the game of checkers. From 1999 to 2001, Fogel and Chellapilla published papers showing how a convolutional neural network could learn to play checker using co-evolution. The learning process did not use prior human professional games, but rather focused on a minimal set of information contained in the checkerboard: the location and type of pieces, and the difference in number of pieces between the two sides. Ultimately, the program (Blondie24) was tested on 165 games against players and ranked in the highest 0.4%. It also earned a win against the program Chinook at its "expert" level of play.
Go
CNNs have been used in computer Go. In December 2014, Clark and Storkey published a paper showing that a CNN trained by supervised learning from a database of human professional games could outperform GNU Go and win some games against Monte Carlo tree search Fuego 1.1 in a fraction of the time it took Fuego to play. Later it was announced that a large 12-layer convolutional neural network had correctly predicted the professional move in 55% of positions, equalling the accuracy of a 6 dan human player. When the trained convolutional network was used directly to play games of Go, without any search, it beat the traditional search program GNU Go in 97% of games, and matched the performance of the Monte Carlo tree search program Fuego simulating ten thousand playouts (about a million positions) per move.
A couple of CNNs for choosing moves to try ("policy network") and evaluating positions ("value network") driving MCTS were used by AlphaGo, the first to beat the best human player at the time.
Time series forecasting
Recurrent neural networks are generally considered the best neural network architectures for time series forecasting (and sequence modeling in general), but recent studies show that convolutional networks can perform comparably or even better. Dilated convolutions might enable one-dimensional convolutional neural networks to effectively learn time series dependences. Convolutions can be implemented more efficiently than RNN-based solutions, and they do not suffer from vanishing (or exploding) gradients. Convolutional networks can provide an improved forecasting performance when there are multiple similar time series to learn from. CNNs can also be applied to further tasks in time series analysis (e.g., time series classification or quantile forecasting).
Cultural heritage and 3D-datasets
As archaeological findings such as clay tablets with cuneiform writing are increasingly acquired using 3D scanners, benchmark datasets are becoming available, including HeiCuBeDa providing almost 2000 normalized 2-D and 3-D datasets prepared with the GigaMesh Software Framework. So curvature-based measures are used in conjunction with geometric neural networks (GNNs), e.g. for period classification of those clay tablets being among the oldest documents of human history.
Fine-tuning
For many applications, training data is not very available. Convolutional neural networks usually require a large amount of training data in order to avoid overfitting. A common technique is to train the network on a larger data set from a related domain. Once the network parameters have converged an additional training step is performed using the in-domain data to fine-tune the network weights, this is known as transfer learning. Furthermore, this technique allows convolutional network architectures to successfully be applied to problems with tiny training sets.
Human interpretable explanations
End-to-end training and prediction are common practice in computer vision. However, human interpretable explanations are required for critical systems such as a self-driving cars. With recent advances in visual salience, spatial attention, and temporal attention, the most critical spatial regions/temporal instants could be visualized to justify the CNN predictions.
Related architectures
Deep Q-networks
A deep Q-network (DQN) is a type of deep learning model that combines a deep neural network with Q-learning, a form of reinforcement learning. Unlike earlier reinforcement learning agents, DQNs that utilize CNNs can learn directly from high-dimensional sensory inputs via reinforcement learning.
Preliminary results were presented in 2014, with an accompanying paper in February 2015. The research described an application to Atari 2600 gaming. Other deep reinforcement learning models preceded it.
Deep belief networks
Convolutional deep belief networks (CDBN) have structure very similar to convolutional neural networks and are trained similarly to deep belief networks. Therefore, they exploit the 2D structure of images, like CNNs do, and make use of pre-training like deep belief networks. They provide a generic structure that can be used in many image and signal processing tasks. Benchmark results on standard image datasets like CIFAR have been obtained using CDBNs.
Neural abstraction pyramid
The feed-forward architecture of convolutional neural networks was extended in the neural abstraction pyramid by lateral and feedback connections. The resulting recurrent convolutional network allows for the flexible incorporation of contextual information to iteratively resolve local ambiguities. In contrast to previous models, image-like outputs at the highest resolution were generated, e.g., for semantic segmentation, image reconstruction, and object localization tasks.
Notable libraries
Caffe: A library for convolutional neural networks. Created by the Berkeley Vision and Learning Center (BVLC). It supports both CPU and GPU. Developed in C++, and has Python and MATLAB wrappers.
Deeplearning4j: Deep learning in Java and Scala on multi-GPU-enabled Spark. A general-purpose deep learning library for the JVM production stack running on a C++ scientific computing engine. Allows the creation of custom layers. Integrates with Hadoop and Kafka.
Dlib: A toolkit for making real world machine learning and data analysis applications in C++.
Microsoft Cognitive Toolkit: A deep learning toolkit written by Microsoft with several unique features enhancing scalability over multiple nodes. It supports full-fledged interfaces for training in C++ and Python and with additional support for model inference in C# and Java.
TensorFlow: Apache 2.0-licensed Theano-like library with support for CPU, GPU, Google's proprietary tensor processing unit (TPU), and mobile devices.
Theano: The reference deep-learning library for Python with an API largely compatible with the popular NumPy library. Allows user to write symbolic mathematical expressions, then automatically generates their derivatives, saving the user from having to code gradients or backpropagation. These symbolic expressions are automatically compiled to CUDA code for a fast, on-the-GPU implementation.
Torch: A scientific computing framework with wide support for machine learning algorithms, written in C and Lua.
See also
Attention (machine learning)
Convolution
Deep learning
Natural-language processing
Neocognitron
Scale-invariant feature transform
Time delay neural network
Vision processing unit
Topological deep learning
Notes
References
External links
CS231n: Convolutional Neural Networks for Visual Recognition — Andrej Karpathy's Stanford computer science course on CNNs in computer vision
vdumoulin/conv_arithmetic: A technical report on convolution arithmetic in the context of deep learning. Animations of convolutions.
Neural network architectures
Computer vision
Computational neuroscience | Convolutional neural network | [
"Engineering"
] | 12,456 | [
"Artificial intelligence engineering",
"Packaging machinery",
"Computer vision"
] |
40,410,426 | https://en.wikipedia.org/wiki/Cohn-Vossen%27s%20inequality | In differential geometry, Cohn-Vossen's inequality, named after Stefan Cohn-Vossen, relates the integral of Gaussian curvature of a non-compact surface to the Euler characteristic. It is akin to the Gauss–Bonnet theorem for a compact surface.
A divergent path within a Riemannian manifold is a smooth curve in the manifold that is not contained within any compact subset of the manifold. A complete manifold is one in which every divergent path has infinite length with respect to the Riemannian metric on the manifold. Cohn-Vossen's inequality states that in every complete Riemannian 2-manifold S with finite total curvature and finite Euler characteristic, we have
where K is the Gaussian curvature, dA is the element of area, and χ is the Euler characteristic.
Examples
If S is a compact surface (without boundary), then the inequality is an equality by the usual Gauss–Bonnet theorem for compact manifolds.
If S has a boundary, then the Gauss–Bonnet theorem gives
where is the geodesic curvature of the boundary, and its integral the total curvature which is necessarily positive for a boundary curve, and the inequality is strict. (A similar result holds when the boundary of S is piecewise smooth.)
If S is the plane R2, then the curvature of S is zero, and χ(S) = 1, so the inequality is strict: 0 < 2.
Notes and references
External links
Gauss–Bonnet theorem, in the Encyclopedia of Mathematics, including a brief account of Cohn-Vossen's inequality
Theorems in differential geometry
Inequalities | Cohn-Vossen's inequality | [
"Mathematics"
] | 336 | [
"Theorems in differential geometry",
"Mathematical theorems",
"Binary relations",
"Mathematical relations",
"Inequalities (mathematics)",
"Theorems in geometry",
"Mathematical problems"
] |
40,411,044 | https://en.wikipedia.org/wiki/Molybdenum%28III%29%20chloride | Molybdenum(III) chloride is the inorganic compound with the formula MoCl3. It forms purple crystals.
Synthesis and structure
Molybdenum(III) chloride is synthesized by the reduction of molybdenum(V) chloride with hydrogen. A higher yield is produced by the reduction of pure molybdenum(V) chloride with anhydrous tin(II) chloride as the reducing agent.
Molybdenum trichloride exists as two polymorphs: alpha (α) and beta (β). The alpha structure is similar to that of aluminum chloride (AlCl3). In this structure, molybdenum has octahedral coordination geometry and exhibits cubic close-packing in its crystalline structure. The beta structure, however, exhibits hexagonal close packing.
Ether complexes
Molybdenum trichloride gives a ether complexes MoCl3(thf)3 and MoCl3(Et2O)3. They are beige, paramagnetic solids. Both feature octahedral Mo centers. The diethyl ether complex is synthesized by reducing a Et2O solution of MoCl5 with tin powder. Older procedures involve stepwise reduction involving isolation of the Mo(IV)-thf complex.
Hexa(tert-butoxy)dimolybdenum(III) is prepared by the salt metathesis reaction from MoCl3(thf)3:
2 MoCl3(thf)3 + 6 LiOBu-t → Mo2(OBu-t)6 + 6 LiCl + 6 thf
References
Chlorides
Molybdenum halides
Molybdenum(III) compounds | Molybdenum(III) chloride | [
"Chemistry"
] | 358 | [
"Chlorides",
"Inorganic compounds",
"Salts"
] |
40,413,995 | https://en.wikipedia.org/wiki/Mallory%27s%20trichrome%20stain | Mallory's trichrome stain also called Mallory's Triple Stain is a stain utilized in histology to aid in revealing different macromolecules that make up the cell. It uses the three stains: aniline blue, acid fuchsin, and orange G. As a result, this staining technique can reveal collagen, ordinary cytoplasm, and red blood cells. It is used in examining the collagen of connective tissue.
For tissues that are not directly acidic or basic, it can be difficult to use only one stain to reveal the necessary structures of interest. A combination of the three different stains in precise amounts applied in the correct order reveals the details selectively. This is the result of more than just electrostatic interactions of stain with the tissue and the stain not being washed out after each step. Collectively the stains complement one another.
The staining technique was first published in 1900 by Frank Burr Mallory, then a histologist at Harvard University Medical School. Many variants of the method exist to simplify or speed processing or to stain other materials. Mallory's and other polychrome stains developed in the early 20th century led to Papanicolaou stain and other popular polychrome staining methods.
The primary application when the stain was introduced was differentiation of structures in connective tissue, and this remains its most common use. Some work however has indicated the stain can highlight differential RNA synthesis. This has been used in identifying ectopic endometrial tissue.
References
Histology
Staining dyes | Mallory's trichrome stain | [
"Chemistry"
] | 312 | [
"Histology",
"Microscopy"
] |
40,417,621 | https://en.wikipedia.org/wiki/Scout%20%28autonomous%20boat%29 | Scout is an autonomous robotic boat designed to complete the first autonomous transatlantic crossing. The project was started by Tiverton students Dylan Rodriguez and Max Kramers in 2010 with the goal of creating an autonomous craft to make the journey from Rhode Island to Sanlucar de Barrameda, Spain. After several iterations, Scout's first transatlantic attempt was launched from Sakonnet Point on June 29, 2013, but unfavourable weather conditions forced the team to recover the craft the same day. A second launch was made July 4, 2013 but after two days a technical failure forced another recovery effort and a redesign of parts of the vessel. The third attempt was launched during the early morning of August 24, 2013 and is currently in progress but has already earned the record for distance of an unmanned Atlantic naval voyage.
Technical specifications
Scout is based on a custom carbon fiber hull with a Divinycell foam core and measures 12.8 feet long, and 25 inches wide and weighs 160 pounds. The vessel design incorporates a bulb keel to right the craft in heavy seas, and propulsion is provided by an electric trolling motor powered by a bank of solar-charged lithium iron phosphate batteries. On-board control and navigation is provided by two Arduino microcontrollers and a GPS receiver, and telemetry data is sent back to the team using the Iridium satellite constellation and provided live on the World Wide Web.
References
External links
Scout Transatlantic
Live Tracking
Unmanned surface vehicles of the United States
Robotics | Scout (autonomous boat) | [
"Engineering"
] | 303 | [
"Robotics",
"Automation"
] |
40,420,141 | https://en.wikipedia.org/wiki/C7H12O3 | {{DISPLAYTITLE:C7H12O3}}
The molecular formula C7H12O3 (molar mass: 144.17 g/mol, exact mass: 144.0786 u) may refer to:
Botryodiplodin
Ethyl levulinate
Tetrahydrofurfuryl acetate
Molecular formulas | C7H12O3 | [
"Physics",
"Chemistry"
] | 73 | [
"Molecules",
"Set index articles on molecular formulas",
"Isomerism",
"Molecular formulas",
"Matter"
] |
35,200,138 | https://en.wikipedia.org/wiki/Nouvelle%20tendance | Nouvelle Tendance (New Tendency) was an art movement founded in Yugoslavia in 1961. The "theoretician" of the group was Croatian art critic Matko Meštrović. The other original founders of Nouvelle Tendance were Brazilian painter Almir Mavignier, and Božo Bek, the Croatian director of the Museum of Contemporary Art, Zagreb.
Overview
The Nouvelle Tendance was more accurately a reflection of common approaches used by artists in a variety of simultaneous movements worldwide, such as Concrete art, kinetic art, and Op Art. The main consideration of the movement has been described as the "problem of movement as conveyed through repetition". The "sensation of displeasure" is provoked in some Nouvelle Tendance works, to "stimulate a more active field of vision" and interest the spectator in an "auto-creative process".
Distinct self-identified groups of artists who became associated with Nouvelle Tendance included GRAV, Gruppo T, Gruppo N, and Zero. The movement attracted artists from France, Germany, Italy, the Netherlands, and Spain. International artists participated in a series of exhibits at European galleries.
This group should not be confused with an early 20th century circle of Paris-based artists who operated briefly under the name "Tendances nouvelles" and held an exhibition in 1904; founding members included Alice Dannenberg and Martha Stettler.
See also
Kinetic art
Op art
Concrete art
Sound art
Sound installation
References
Visual arts genres
Modern art
Types of sculpture
Motion (physics)
Contemporary art movements | Nouvelle tendance | [
"Physics"
] | 308 | [
"Physical phenomena",
"Motion (physics)",
"Space",
"Mechanics",
"Spacetime"
] |
35,201,597 | https://en.wikipedia.org/wiki/Force%20matching | Force matching is a research method consisting of test subjects attempting to produce a set forces that are equal to a set of more reliable reference forces.
Types
Biomechanical
A subject’s maximum voluntary contraction (MVC) is recorded and used to normalize both reference forces and results between subjects. During the test subjects are assisted in producing a reference force using various types of feedback (static weight or visual display of force generated). This is followed by an attempt of the subject to generate the reference force without assistance. The duration for both reference and matching tasks is usually four seconds. Results are taken as a mean value of force generated over a time interval set by the researcher. Time intervals are generally one second long and near the end of the attempt. Reference forces are typically set as a percentage of a subject’s MVC while error is typically reported as a percentage of a subject’s MVC.
Atomic
It is one of the effective research method to obtain realistic classical interatomic potential or force field for molecular dynamics simulation with high degree of transferability for systems which the first principles or ab initio method is capable of treating. This method is based on fitting the forces on individual atoms in a number of reference structures, cohesive energies and stresses on unit cell obtained from first principles calculation with those obtained from classical interatomic potential. The target of the computational fitting is to determine unknown coefficients in classical interatomic potential function. This method is developed by F. Ercolessi, and J. B. Adams during 1992 and 1993 at Department of Material Science and Engineering at the University of Illinois Urbana-Champaign. The enormous number of reference structures, which can reach several thousand values, makes it possible to fit large number of parameters needed for potential in binary and ternary systems.
For Lennard-Jones potential:
where ε is the depth of the potential well, σ is the finite distance at which the inter-particle potential is zero, r is the distance
between the particles. These two unknown parameters can be fitted to reproduce experimental data or accurate data obtained from
first principle calculations. Differentiating the L-J potential with respect to r gives an expression for the net inter-molecular force between 2 molecules. This inter-molecular force may be attractive or repulsive, depending on the value of r. When r is very small, the molecules repel each other. In force matching method the forces from classical potential
are compared with reference force calculated from ab initio method to determine and .
Applications
Biomechanical force matching has been used by researchers to describe the accuracy of muscle contractions under various conditions. It has been observed that the thumb is more accurate in force matching than fingers are. Impairment of the extensor pollicis longus has not produced a decrease in force matching accuracy of the flexor pollicis longus.
Notes
References
Research methods
Force | Force matching | [
"Physics",
"Mathematics"
] | 579 | [
"Force",
"Physical quantities",
"Quantity",
"Mass",
"Classical mechanics",
"Wikipedia categories named after physical quantities",
"Matter"
] |
22,013,888 | https://en.wikipedia.org/wiki/Cerrolow%20136 | Cerrolow 136, is a fusible alloy that becomes liquid at approximately 58 °C (136 °F). It is a eutectic alloy of bismuth, lead, indium, and tin, with the following percentages by weight: 49% Bi, 18% Pb, 21% In, 12% Sn.
It is also known as ChipQuik desoldering alloy. and Lens Alloy 136, used for mounting lenses and other optical components for grinding. Used for mounting small delicate oddly-shaped components for machining. It slightly expands on cooling, later shows slight shrinkage in couple hours afterwards. Used as a solder in low-temperature physics.
The Young's Modulus is approximately 2.5×106 psi (17.24 GPa)
Similar metals
References
Fusible alloys | Cerrolow 136 | [
"Chemistry",
"Materials_science"
] | 166 | [
"Metallurgy",
"Alloys",
"Fusible alloys"
] |
22,016,953 | https://en.wikipedia.org/wiki/J1%20J2%20model | The J1–J2 model is a quantum spin model like the Heisenberg model but also includes a term for the interaction between next-nearest neighbor spins.
Hamiltonian
In this model, the term represents the usual nearest-neighbor interaction as seen in the Heisenberg model, and represents the exchange interaction to the next nearest-neighbor.
See also
Spin model
Heisenberg model (quantum)
Hubbard model
t-J model
Majumdar–Ghosh model
References
Spin models
Quantum magnetism | J1 J2 model | [
"Physics",
"Materials_science"
] | 101 | [
"Spin models",
"Quantum mechanics",
"Quantum magnetism",
"Condensed matter physics",
"Statistical mechanics"
] |
37,650,848 | https://en.wikipedia.org/wiki/Automated%20parking%20system | An automated (car) parking system (APS) is a mechanical system designed to minimize the area and/or volume required for parking cars. Like a multi-story parking garage, an APS provides parking for cars on multiple levels stacked vertically to maximize the number of parking spaces while minimizing land usage. The APS, however, utilizes a mechanical system to transport cars to and from parking spaces (rather than the driver) in order to eliminate much of the space wasted in a multi-story parking garage. While a multi-story parking garage is similar to multiple parking lots stacked vertically, an APS is more similar to an automated storage and retrieval system for cars. Parking systems are generally powered by electric motors or hydraulic pumps that move vehicles into a storage position.The paternoster (shown animated at the right) is an example of one of the earliest and most common types of APS.
APS are also generically known by a variety of other names, including:automated parking facility (APF), automated vehicle storage and retrieval system (AVSRS), car parking system, mechanical parking, and robotic parking garage.
History
The concept for the automated parking system was and is driven by two factors: a need for parking spaces and a scarcity of available land. The earliest use of an APS was in Paris, France in 1905 at the Garage Rue de Ponthieu. The APS consisted of a groundbreaking multi-story concrete structure with an internal car elevator to transport cars to upper levels where attendants parked the cars.
In the 1920s, a Ferris wheel-like APS (for cars rather than people) called a paternoster system became popular as it could park eight cars in the ground space normally used for parking two cars. Mechanically simple with a small footprint, the paternoster was easy to use in many places, including inside buildings. At the same time, Kent Automatic Garages was installing APS with capacities exceeding 1,000 cars. The “ferris-wheel,” or paternoster system — was created by the Westinghouse Corporation in 1923 and subsequently built in 1932 on Chicago's Monroe Street. The Nash Motor Company created the first glass-enclosed version of this system for the Chicago Century of Progress Exhibition in 1933
The first driverless parking garage opened in 1951 in Washington, D.C., but was replaced with office space due to increasing land values.
APS saw a spurt of interest in the U.S. in the late 1940s and 1950s with the Bowser, Pigeon Hole and Roto Park systems. In 1957, 74 Bowser, Pigeon Hole systems were installed, and some of these systems remain in operation. However, interest in APS in the U.S. waned due to frequent mechanical problems and long waiting times for patrons to retrieve their cars. In the United Kingdom, the Auto Stacker opened in 1961 in Woolwich, south east London, but proved equally difficult to operate. Interest in APS in the U.S. was renewed in the 1990s, and there were 25 major current and planned APS projects (representing nearly 6,000 parking spaces) in 2012. The first American robotic parking garage opened in 2002 in Hoboken, New Jersey.
While interest in the APS in the U.S. languished until the 1990s, Europe, Asia and Central America had been installing more technically advanced APS since the 1970s. In the early 1990s, nearly 40,000 parking spaces were being built annually using the paternoster APS in Japan. In 2012, there are an estimated 1.6 million APS parking spaces in Japan.
The ever-increasing scarcity of available urban land (urbanization) and increase of the number of cars in use (motorization) have combined with sustainability and other quality-of-life issues to renew interest in APS as alternatives to multi-storey car parks, on-street parking, and parking lots.
Largest systems
The largest Automated Parking Facility in the world is in Al Jahra, Kuwait, and provides 2,314 parking spaces.
The world's fastest Automated Parking System is in Wolfsburg, Germany, with a retrieval time of 1 minute and 44 seconds.
The largest APS in Europe is at Dokk1 in Aarhus, Denmark, and provides 1,000 parking spaces via 20 car lifts.
Space saving
All APS take advantage of a common concept to decrease the area of parking spaces - removing the driver and passengers from the car before it is parked. With either fully automated or semi-automated APS, the car is driven up to an entry point to the APS and the driver and passengers exit the car. The car is then moved automatically or semi-automatically (with some attendant action required) to its parking space.
The space-saving provided by the APS, compared to the multi-story parking garage, is derived primarily from a significant reduction in space not directly related to the parking of the car:
Parking space width and depth (and distances between parking spaces) are dramatically reduced since no allowance need be made for driving the car into the parking space or for the opening of car doors (for drivers and passengers)
No driving lanes or ramps are needed to drive the car to/from the entrance/exit to a parking space
Ceiling height is minimized since there is no pedestrian traffic (drivers and passengers) in the parking area, and
No walkways, stairways or elevators are needed to accommodate pedestrians in the parking area.
With the elimination of ramps, driving lanes, pedestrians and the reduction in ceiling heights, the APS requires substantially less structural material than the multi-story parking garage. Many APS utilize a steel framework (some use thin concrete slabs) rather than the monolithic concrete design of the multi-story parking garage. These factors contribute to an overall volume reduction and further space savings for the APS.
Other considerations
In addition to the space saving, many APS designs provide a number of secondary benefits:
The parked cars and their contents are more secure since there is no public access to parked cars
Minor parking lot damage such as scrapes and dents are eliminated
Drivers and passengers are safer not having to walk through parking lots or garages
Driving around in search of a parking space is eliminated, thereby reducing engine emissions and wasted time
Only minimal ventilation and lighting systems are needed
Handicap access is improved
The volume and visual impact of the parking structure is minimized
Shorter construction time
Problems
There have been a number of problems with robotic parking systems, particularly in the United States. The systems work well in balanced throughput situations like shopping malls and train stations, but they are unsuited to high peak volume applications like rush hour usage or stadiums and they suffer from technical problems. Further, parkers not familiar with the system may cause problems, for example by failing to push the button to alert a fully automated system to the presence of a car to be parked.
Fully automated vs semi-automated
Fully automated parking systems operate much like robotic valet parking. The driver drives the car into an APS entry (transfer) area. The driver and all passengers exit the car. The driver uses an automated terminal nearby for payment and receipt of a ticket. When driver and passengers have left the entry area, the mechanical system lifts the car and transports it to a pre-determined parking space in the system. More sophisticated fully automated APS will obtain the dimensions of cars on entry in order to place them in the smallest available parking space.
The driver retrieves a car by inserting a ticket or code into an automated terminal. The APS lifts the car from its parking space and delivers it to an exit area. Most often, the retrieved car has been oriented to eliminate the need for the driver to back out.
Fully automated APS theoretically eliminate the need for parking attendants.
Semi-automated APS also use a mechanical system of some type to move a car to its parking space, however putting the car into and/or the operation of the system requires some action by an attendant or the driver.
The choice between fully and semi-automated APS is often a matter of space and cost, however large capacity (> 100 cars) tend to be fully automated.
Applications
By virtue of their relatively smaller volume and mechanized parking systems, APS are often used in locations where a multi-story parking garage would be too large, too costly or impractical. Examples of such applications include, under or inside existing or new structures, between existing structures and in irregularly shaped areas.
APS can also be applied in situations similar to multi-storey parking garages such as freestanding above ground, under buildings above grade and under buildings below grade.
Costs
The direct comparison of costs between an APS and a multi-story parking garage can be complicated by many variables such as capacity, land costs, area shape, number and location of entrances and exits, land usage, local codes and regulations, parking fees, location, and aesthetic and environmental requirements.
Following is a comparison of building costs for generic APS and multi-story parking garages:
The comparison above is for building costs only. Not included, for example, is the cost of land or the opportunity cost of the use of the land (i.e. the value of the additional space made available by the smaller size of the APS). As evidence of the complexities of comparing the costs for APS and multi-story parking garages, the same author presents an actual case study as follows:
In this case study, the APS also provides roughly of additional open space compared to the multi-story parking garage which provides no open space and requires minimum setbacks be utilized. Other references also indicate that the cost comparison between APS and multi-story parking garages is highly dependent on the application and the detailed design.
See also
Automated storage and retrieval system
Sustainability
References
Parking
Automation | Automated parking system | [
"Engineering"
] | 1,992 | [
"Control engineering",
"Automation"
] |
37,656,184 | https://en.wikipedia.org/wiki/Laminin%20111 | Laminin–111 (also "laminin–1") is a protein of the type known as laminin isoforms. It was among the first of the laminin isoforms to be discovered. The "111" identifies the isoform's chain composition of α1β1γ1. This protein plays an important role in embryonic development. Injections of this substance are used in treatment for Duchenne muscular dystrophy, and its cellular action may potentially become a focus of study in cancer research.
Distribution
The distribution of the different laminin isoforms is tissue-specific. Laminin–111 is predominantly expressed in the embryonic epithelium, but can also be found in some adult epithelium such as the kidney, liver, testis, ovaries, and brain blood vessels. Different levels of expression of α chains have a large influence on the differential expression of laminin, thereby determining the isoform produced. From studying a mouse model, it was found that transcription factors present in the parietal endoderm regulate the expression of the α1 and large amounts of laminin-111 are produced.
Functions
The synthesized laminin–111 formed in an embryo contributes to the formation of Reichert’s membrane, a thick extra-embryonic basement membrane. When the laminin α1 chain is deficient in an organism, an embryo dies, likely as a result of a defective Reichert’s membrane due to a lack of laminin–111 being produced. Laminin-111 has been identified as a crucial molecule for development of the embryo as shown by the consequences that occur when laminin-111 is lacking.
Laminin-111 is expressed very early on in development and is present in the blastocyst. When various parts of the trimer chains are knocked out by mutations, devastating consequences occur in the embryo. If the β1 or γ1 chains of laminin-111 are absent the basement membrane fails to form. Without a basement membrane cells have nowhere to attach and all dependent activities such as cell migration and epithelial formation can no longer occur. The self-assembly and tight network formation by laminin-111 are essential for holding the basement membrane together.
Although it is expressed abundantly during the early embryonic stage, laminin-111 is mostly absent in adults. The injection of laminin-111, however, helps with Duchenne muscular dystrophy, a neuromuscular disease in which the connection between the extracellular matrix and cell cytoskeleton is lost. Increased levels of laminin-111 triggered an increase in the expression of α7-integrin receptor and this prevented onset of the disease. Additionally, the presence of laminin-111 increased muscle strength and protected it from injury. When injected with myoblast transplants, laminin–111 decreased degeneration and inflammatory reactions and increased the success of the transplantation. The experiments utilizing laminin–111 as a source of therapy for Duchenne muscular dystrophy suggest that it has protective qualities in addition to its association with muscle tissue.
Mechanisms of action
Cell adhesion
In cell adhesion laminin-111 and other isoforms are important proteins that anchor cells to the extracellular matrix (ECM). The linkage between cells and the ECM is formed by binding cell surface receptors to one end of the laminin α chain and binding ECM components to another region of the laminin. Globular domains (G-Domain) of the α chain are the regions on laminin-111 that allow the binding of integrins, glycoproteins, sulfated glycolipids and dystroglycan.
Cell signaling
Besides anchoring cells to the ECM, laminins are also involved in the signalling of cells and other components of the ECM. Even though there is not a general mechanism that applies to all laminins in signalling, there are some common pathways that can be seen in more than one isoform of laminin. For example, the PI3K/AKT pathway is used by laminin-111 (promotes cell-survival), 511 (prevents apoptosis with laminin 521), and 521 (stabilizes pluripotency of human embryonic stem cells). The pathway begins with the adhesion of the cell to the ECM for activation of the lipid-associated PI3K. Once PI3K is activated, it will localize AKT that is in the cytoplasm to the cell membrane where AKT is then phosphorylated to promote cell survival.
Neurite outgrowth
When α chains of laminin-111 bind to cell surface receptors integrins α1β1, α3β1, α4β1, α6β1 and Cdc42 GTPase are activated. The activated GTPase then activates Cdc42 which further activates c-Jun kinases and phosphorylation of Jun. Activation of c-Jun kinases leads to high levels of c-Jun expression which results in neurite outgrowth. The synthesis of Nitric Oxides resides somewhere in the pathway and is yet to be determined. Weston et al. (2000) proposed that the synthesis of Nitric Oxide may be upstream to the activation of Cdc42. Nonetheless, Nitric Oxide synthesis is shown to be an important element in laminin-mediated neurite outgrowth.
Future applications
Dynamic reciprocity theory
The dynamic reciprocity theory states that a cell’s fate depends on the exchange of chemical signals between the extracellular matrix and the nucleus of the cell. Focussing on connections between laminin-111 and other proteins involved in cell-to-cell communication could spark further research that may help to further our current understanding of cancer and how to slow down or stop its process.
Actin plays a role in nuclear activity which is an important process with regard to cell signalling influencing cell differentiation and replication. It has been suggested that actin interactions directly influence gene transcription as it interacts with chromatin remodeling complexes as well as RNA polymerases I, II and III. However, the exact role that actin plays in transcription has not yet been determined.
Implications for cancer research
A group of distinguished scientists from the U.S. Department of Energy’s (DOE) Lawrence Berkeley National Laboratory (Berkeley Lab) did a recent study on how laminin-111 interacts with the cytoplasmic protein, actin. Their study gave the following conclusions:
The biological process in which a cell ceases to continue growing and dividing is called quiescence (the opposite of cancer). ECM laminin-111 sends chemical signals that promotes adhesion of a cell and its ECM. Although the mechanism is unknown, these signals have also been linked to cell quiescence. Adding laminin-111 to breast epithelial cells leads to quiescence by altering nuclear actin. High levels of laminin-111 deplete nuclear actin which induces quiescence of cells. However, when an isoform of actin, that cannot exit a cell’s nucleus, is active, cells continue to grow and divide even when laminin levels are high. ECM laminin-111 levels in a normal breast cell are significantly higher than laminin-111 levels in tissues of cancerous breast tissue. Simply increasing laminin levels in the ECM of cancerous breast cells is not enough to lead to quiescence. Therefore, it is implied that there are multiple factors working together influencing cell-to-cell communication. How laminin-111 and nuclear actin communicate is one of these factors. Laminin-111 could be the physiological regulator of nuclear actin which would suggest that depleting nuclear actin could be a key to achieving cell quiescence and returning to homeostatic operating conditions. Decreased expression of laminin-111 and the growth-inhibitory signals that it produces in malignant myoepithelial cells begs for further investigation with regard to cancer research. Therefore, further exploration of laminin-111 and nuclear actin interaction could be a target for future experimental therapeutic investigations.
References
Carbohydrate chemistry | Laminin 111 | [
"Chemistry"
] | 1,720 | [
"Carbohydrate chemistry",
"nan",
"Chemical synthesis",
"Glycoproteins",
"Glycobiology"
] |
37,656,368 | https://en.wikipedia.org/wiki/Volcanic%20impacts%20on%20the%20oceans | Explosive volcanic eruptions affect the global climate in several ways.
Lowering sea surface temperature
One main impact of volcanoes is the injection of sulfur-bearing gases into the stratosphere, which oxidize to form sulfate aerosols. Stratospheric sulfur aerosols spread around the globe by the atmospheric circulation, producing surface cooling by scattering solar radiation back to space. This cooling effect on the ocean surface usually lasts for several years as the lifetime of sulfate aerosols is about 2–3 years. However, in the subsurface ocean the cooling signal may persist for a longer time and may have impacts on some decadal variabilities, such as the Atlantic meridional overturning circulation (AMOC).
Volcanic aerosols from huge volcanoes (VEI>=5) directly reduce global mean sea surface temperature (SST) by approximately 0.2-0.3 °C, milder than global total surface temperature drop, which is ~0.3 to 0.5 °C, according to both global temperature records and model simulations. It usually takes several years to be back to normal.
Decreasing ocean heat content
The volcanic cooling signals in ocean heat content can persist for much longer time (decadal or mutil-decadal time scale), far beyond the duration of volcanic forcing.
Several studies have revealed that Krakatau’s effect in the heat content can be as long as one-century. Relaxation time of the effects of recent volcanoes is generally shorter than those before the 1950s. For example, the recovery time of ocean heat content of Pinatubo, which caused comparable radiative forcing to Krakatau, seems to be much shorter. This is because Pinatubo happened under a warm and non-stationary background with increasing greenhouse gas forcing. However, its signal still could penetrate down to ~1000 m deep.
A 2022 study on environmental impacts of volcanic eruptions showed that in the eastern equatorial of the pacific, after the volcano erupts, some low-latitude volcano trends to warmer. But some highlatitude volcanoes tend to be colder.
Altering sea level
As thermal expansion is a key factor in sea level variability, decreased heat content should result in a reduction in global mean sea level on a decadal time scale. However, Grinsted [2007] argued that a significant sea level rise is the first direct response to the volcanic eruption, and after that sea level becomes to drop. One possible explanation for this phenomenon is the imbalance of ocean mass fluxes. After a volcanic eruption, evaporation over ocean will lower, because it is largely determined by the ocean surface temperature change. The quick response of evaporation to the surface cooling and the delayed response of river runoff to the associated lower precipitation lead to an increased sea level. About 1~2 years later, river discharge becomes less due to the reduced precipitation and less sea ice melting, which cause sea level to drop.
Ocean oxygen and carbon levels
In 2023, several scientists studied the eruption of Mt Pinatubo in June 1991 and discovered that it led to increases in the ocean oxygen and carbon concentrations that persisted for many years.
Enhancing AMOC
Results from a number of modeling studies suggest that the Atlantic Meridional Overturning Circulation (AMOC) is enhanced by volcanic activity. The deepwater formation at the northern end of the Atlantic Ocean allows SST anomalies to be subducted into the deep ocean efficiently because the rate of overturning is altered by changes in salinity. The decreasing summer-time ice melting and precipitation due to the volcano cooling enhance the salinity near the Greenland Sea, and further reduces static stability, which means more surface water sinks into the deep ocean. The studies of Stenchikov et al. (2009) and Iwi (2012) suggest that both Krakatau and Pinatubo may have strengthened the overturning circulation. And the increase in AMOC seems to be strongest at about one decade after the volcano eruption, with a magnitude of about one sverdrup for Krakatau and Pinatubo.
References
Volcanic events
Oceanography | Volcanic impacts on the oceans | [
"Physics",
"Environmental_science"
] | 840 | [
"Oceanography",
"Hydrology",
"Applied and interdisciplinary physics"
] |
37,657,202 | https://en.wikipedia.org/wiki/Planetary%20equilibrium%20temperature | The planetary equilibrium temperature is a theoretical temperature that a planet would be if it were in radiative equilibrium, typically under the assumption that it radiates as a black body being heated only by its parent star. In this model, the presence or absence of an atmosphere (and therefore any greenhouse effect) is irrelevant, as the equilibrium temperature is calculated purely from a balance with incident stellar energy.
Other authors use different names for this concept, such as equivalent blackbody temperature of a planet. The effective radiation emission temperature is a related concept, but focuses on the actual power radiated rather than on the power being received, and so may have a different value if the planet has an internal energy source or when the planet is not in radiative equilibrium.
Planetary equilibrium temperature differs from the global mean temperature and surface air temperature, which are measured observationally by satellites or surface-based instruments, and may be warmer than the equilibrium temperature due to the greenhouse effect.
Calculation of equilibrium temperature
Consider a planet orbiting its host star. The star emits radiation isotropically, and some fraction of this radiation reaches the planet. The amount of radiation arriving at the planet is referred to as the incident solar radiation, . The planet has an albedo that depends on the characteristics of its surface and atmosphere, and therefore only absorbs a fraction of radiation. The planet absorbs the radiation that isn't reflected by the albedo, and heats up. One may assume that the planet radiates energy like a blackbody at some temperature according to the Stefan–Boltzmann law. Thermal equilibrium exists when the power supplied by the star is equal to the power emitted by the planet. The temperature at which this balance occurs is the planetary equilibrium temperature.
Derivation
The solar flux absorbed by the planet from the star is equal to the flux emitted by the planet:
Assuming a fraction of the incident sunlight is reflected according to the planet's Bond albedo, :
where represents the area- and time-averaged incident solar flux, and may be expressed as:
The factor of 1/4 in the above formula comes from the fact that only a single hemisphere is lit at any moment in time (creates a factor of 1/2), and from integrating over angles of incident sunlight on the lit hemisphere (creating another factor of 1/2).
Assuming the planet radiates as a blackbody according to the Stefan–Boltzmann law at some equilibrium temperature , a balance of the absorbed and outgoing fluxes produces:
where is the Stefan-Boltzmann constant.
Rearranging the above equation to find the equilibrium temperature leads to:
where is the luminosity of the Sun ( W), and the distance between the planet and the Sun, then :
(with in metres), or :
(with in million kilometres).
Calculation for extrasolar planets
For a planet around another star, (the incident stellar flux on the planet) is not a readily measurable quantity. To find the equilibrium temperature of such a planet, it may be useful to approximate the host star's radiation as a blackbody as well, such that:
The luminosity () of the star, which can be measured from observations of the star's apparent brightness, can then be written as:
where the flux has been multiplied by the surface area of the star.
To find the incident stellar flux on the planet, , at some orbital distance from the star, , one can divide by the surface area of a sphere with radius :
Plugging this into the general equation for planetary equilibrium temperature gives:
If the luminosity of the star is known from photometric observations, the other remaining variables that must be determined are the Bond albedo and orbital distance of the planet. Bond albedos of exoplanets can be constrained by flux measurements of transiting exoplanets, and may in future be obtainable from direct imaging of exoplanets and a conversion from geometric albedo. Orbital properties of the planet such as the orbital distance can be measured through radial velocity and transit period measurements.
Alternatively, the planetary equilibrium may be written in terms of the temperature and radius of the star:
Caveats
The equilibrium temperature is neither an upper nor lower bound on actual temperatures on a planet. There are several reasons why measured temperatures deviate from predicted equilibrium temperatures.
Greenhouse effect
In the greenhouse effect, long wave radiation emitted by a planet is absorbed by certain gases in the atmosphere, reducing longwave emissions to space. Planets with substantial greenhouse atmospheres emit more longwave radiation at the surface than what reaches space. Consequently, such planets have surface temperatures higher than their effective radiation emission temperature. For example, Venus has an effective temperature of approximately , but a surface temperature of . Similarly, Earth has an effective temperature of , but a surface temperature of about due to the greenhouse effect in our lower atmosphere. The surface temperatures of such planets are more accurately estimated by modeling thermal radiation transport through the atmosphere.
Airless bodies
On airless bodies, the lack of any significant greenhouse effect allows equilibrium temperatures to approach mean surface temperatures, as on Mars, where the equilibrium temperature is and the mean surface temperature of emission is . There are large variations in surface temperature over space and time on airless or near-airless bodies like Mars, which has daily surface temperature variations of 50–60 K. Because of a relative lack of air to transport or retain heat, significant variations in temperature develop. Assuming the planet radiates as a blackbody (i.e. according to the Stefan-Boltzmann law), temperature variations propagate into emission variations, this time to the power of 4. This is significant because our understanding of planetary temperatures comes not from direct measurement of the temperatures, but from measurements of the fluxes. Consequently, in order to derive a meaningful mean surface temperature on an airless body (to compare with an equilibrium temperature), a global average surface emission flux is considered, and then an 'effective temperature of emission' that would produce such a flux is calculated. The same process would be necessary when considering the surface temperature of the Moon, which has an equilibrium temperature of , but can have temperatures of in the daytime and at night. Again, these temperature variations result from poor heat transport and retention in the absence of an atmosphere.
Internal energy fluxes
Orbiting bodies can also be heated by tidal heating, geothermal energy which is driven by radioactive decay in the core of the planet, or accretional heating. These internal processes will cause the effective temperature (a blackbody temperature that produces the observed radiation from a planet) to be warmer than the equilibrium temperature (the blackbody temperature that one would expect from solar heating alone).
For example, on Saturn, the effective temperature is approximately 95 K, compared to an equilibrium temperature of about 63 K. This corresponds to a ratio between power emitted and solar power received of ~2.4, indicating a significant internal energy source. Jupiter and Neptune have ratios of power emitted to solar power received of 2.5 and 2.7, respectively.
Close correlation between the effective temperature and equilibrium temperature of Uranus can be taken as evidence that processes producing an internal flux are negligible on Uranus compared to the other giant planets.
Earth has insufficient geothermal heating to significantly affect its global temperature, with geothermal heating supplying only 0.03% of Earth's total energy budget.
See also
Earth's energy budget
Effective temperature
Thermal equilibrium
References
Sources
Thermodynamics
Planetary science
Equations of astronomy | Planetary equilibrium temperature | [
"Physics",
"Chemistry",
"Astronomy",
"Mathematics"
] | 1,510 | [
"Concepts in astronomy",
"Thermodynamics",
"Equations of astronomy",
"Planetary science",
"Astronomical sub-disciplines",
"Dynamical systems"
] |
37,657,545 | https://en.wikipedia.org/wiki/C21H19ClN4O2 | {{DISPLAYTITLE:C21H19ClN4O2}}
The molecular formula C21H19ClN4O2 (molar mass: 394.86 g/mol) may refer to:
SB-242084
Setanaxib
SSR-180,575
Molecular formulas | C21H19ClN4O2 | [
"Physics",
"Chemistry"
] | 67 | [
"Molecules",
"Set index articles on molecular formulas",
"Isomerism",
"Molecular formulas",
"Matter"
] |
37,658,021 | https://en.wikipedia.org/wiki/Null%20model | In mathematics, for example in the study of statistical properties of graphs, a null model is a type of random object that matches one specific object in some of its features, or more generally satisfies a collection of constraints, but which is otherwise taken to be an unbiasedly random structure. The null model is used as a term of comparison, to verify whether the object in question displays some non-trivial features (properties that wouldn't be expected on the basis of chance alone or as a consequence of the constraints), such as community structure in graphs. An appropriate null model behaves in accordance with a reasonable null hypothesis for the behavior of the system under investigation.
One null model of utility in the study of complex networks is that proposed by Newman and Girvan, consisting of a randomized version of an original graph , produced through edges being rewired at random, under the constraint that the expected degree of each vertex matches the degree of the vertex in the original graph.
The null model is the basic concept behind the definition of modularity, a function which evaluates the goodness of partitions of a graph into clusters. In particular, given a graph and a specific community partition (an assignment of a community-index (here taken as an integer from to ) to each vertex in the graph), the modularity measures the difference between the number of links from/to each pair of communities, from that expected in a graph that is completely random in all respects other than the set of degrees of each of the vertices (the degree sequence). In other words, the modularity contrasts the exhibited community structure in with that of a null model, which in this case is the configuration model (the maximally random graph subject to a constraint on the degree of each vertex).
See also
Null hypothesis
References
Graph theory | Null model | [
"Mathematics"
] | 369 | [
"Graph theory stubs",
"Discrete mathematics",
"Graph theory",
"Combinatorics",
"Mathematical relations"
] |
37,658,639 | https://en.wikipedia.org/wiki/Adaptive%20Coloration%20in%20Animals | Adaptive Coloration in Animals is a 500-page textbook about camouflage, warning coloration and mimicry by the Cambridge zoologist Hugh Cott, first published during the Second World War in 1940; the book sold widely and made him famous.
The book's general method is to present a wide range of examples from across the animal kingdom of each type of coloration, including marine invertebrates and fishes as well as terrestrial insects, amphibians, reptiles, birds and mammals. The examples are supported by many of Cott's own drawings, diagrams, and photographs. This essentially descriptive natural history treatment is supplemented with accounts of experiments by Cott and others. The book had few precedents, but to some extent follows (and criticises) Abbott Handerson Thayer's 1909 Concealing-Coloration in the Animal Kingdom.
The book is divided into three parts: concealment, advertisement, and disguise. Part 1, concealment, covers the methods of camouflage, which are colour resemblance, countershading, disruptive coloration, and shadow elimination. The effectiveness of these, arguments for and against them, and experimental evidence, are described. Part 2, advertisement, covers the methods of becoming conspicuous, especially for warning displays in aposematic animals. Examples are chosen from mammals, insects, reptiles and marine animals, and empirical evidence from feeding experiments with toads is presented. Part 3, disguise, covers methods of mimicry that provide camouflage, as when animals resemble leaves or twigs, and markings and displays that help to deflect attack or to deceive predators with deimatic displays. Both Batesian mimicry and Müllerian mimicry are treated as adaptive resemblance, much like camouflage, while a chapter is devoted to the mimicry and behaviour of the cuckoo. The concluding chapter admits that the book's force is cumulative, consisting of many small steps of reasoning, and being a wartime book, compares animal to military camouflage.
Cott's textbook was at once well received, being admired both by zoologists and naturalists and among allied soldiers. Many officers carried a copy of the book with them in the field. Since the war it has formed the basis for experimental investigation of camouflage, while its breadth of coverage and accuracy have ensured that it remains frequently cited in scientific papers.
The book
Approach
Adaptive Coloration in Animals is a 500-page book, in its first edition. It was published by Methuen (in London) and Oxford University Press (in New York) in 1940. It is full of detailed observations of types of camouflage and other uses of colour in animals, and illustrated by the author with clear drawings and photographs. There is a coloured frontispiece showing eight of Cott's paintings of tropical amphibians. The book has 48 monotone plates and several illustrations.
Cott's method is to provide a large number of examples, illustrated with his own drawings or photographs, showing animals from different groups including fish, reptiles, birds and insects, especially butterflies. The examples are chosen to illustrate specific adaptations. For example, the fish Chaetodon capistratus is described as follows:
Cott was well aware that he was publishing in wartime. There are, as Julian Huxley remarks in his 'Introduction', references throughout the book to the human analogues of animal camouflage and concealment. For example, in the section on 'Adaptive Silence', the kestrel is said to "practise dive-bombing attacks", or "after the fashion of a fighter 'plane" to fly down other birds, while "Owls have solved the problem of the silent air-raid"; Cott spends the rest of that paragraph on the "method which has recently been rediscovered and put into practice" of shutting off a bomber's engines and "gliding noiselessly down towards their victims" at Barcelona in the Spanish Civil War. In the concluding chapter, Cott explicitly states "The innumerable visible devices used ... in peacetime and in wartime ... are merely rediscovered ... applications of colour that have already reached a high ... degree of specialization and perfection.. in the animal world", mentioning predator-prey relationships, sexual selection and signalling to rivals. He then compares the "hunting disguises put on ... as a means of approaching, ambushing or alluring game, and the sniping suits, concealed machine-gun posts, and booby traps" with the camouflage of animal predators; and similarly he compares "protective disguises" with the "photographer's hide and the gunner's observation post." In the same section, Cott compares intentionally visible signs with animal warning colours: "The policeman's white gloves have their parallel in the white stripes or spots of nocturnal skunks and carabids. The Automobile Association has adopted a system of coloration [black and yellow] whose copyright belongs by priority to wasps and salamanders."
Structure
The book addresses its subject under three main headings: concealment, advertisement, and disguise.
Part I: Concealment
The methods by which concealment is attained in nature
Cott sets out his view that we have to be re-taught how to see, mentioning Ruskin's "innocence of the eye". He argues that camouflage should, and in animals actually does, use four mechanisms: colour resemblance, obliterative shading (i.e. countershading, the graded shading which conceals self-shadowing of the lower body), disruptive coloration, and shadow elimination.
Chapter 1. General colour resemblance.
Cott gives many examples such as a table of 16 species of green tropical tree-snakes.
Chapter 2. Variable colour resemblance. Caterpillars and pupae (as in Poulton's famous experiment) are coloured to match their environment. Mountain hares change colour in winter; many fish, cephalopods, frogs, and crustacea can change colour rapidly.
Chapter 3. Obliterative shading.
Following the artist and amateur naturalist Abbott Handerson Thayer, Cott explains countershading with diagrams, photographs of models and examples of real animals. He shows how helpful it would be for military camouflage with drawings of gun barrels.
Chapter 4. Disruptive coloration.
Cott argues with diagrams, drawings, photographs and examples that animals are often extremely effectively disruptively patterned. He analyses the component effects of disruption, including "differential blending" and "maximum disruptive contrast". Cott's figure 7 is a set of nine drawings, arranged as a 3x3 table. On the left is an animal's outline in grey tone against a differently coloured background. In the centre, the same animals are now disruptively patterned against the same plain backgrounds. On the right, the disruptively patterned animals are shown against realistic broken backgrounds containing vegetation or rocks. Cott explains
Cott goes on to explain that the right-hand drawing shows the effect "of broken surroundings in further blending and confusing the picture", observing that this is the closest to what is seen in nature. His readers are invited to look first at the right-hand images to gain an idea of the power of "these optical devices" as camouflage, putting off the moment when the animal is actually recognised.
Chapter 5. Coincident disruptive coloration.
Animals such as frogs are patterned so that when they are at rest with legs tucked in, their outline is powerfully disrupted with markings that seem to flow across body and leg boundaries. Eyes too are often hidden in stripes or eye masks.
Chapter 6. Concealment Of the shadow.
Cast shadows give away even well-camouflaged animals. Many animals therefore take care to minimise shadow, by lying down, with flattened bodies, or with fringes. Some hawkmoth caterpillars have false shadow patterns to suggest they are parts of other objects.
The function of concealing coloration in nature
Chapter 7. Concealment in defence, mainly as illustrated by birds.
Cott considers how effective camouflage is as an adaptation, such as in incubation and rest (sleep) in birds. For instance nightjars are nocturnal, and rest, well camouflaged, on the ground during the day.
Chapter 8. Concealment In offence.
Cott describes the care that predators take when approaching prey, minimizing visible movement and scent, the use of cover for ambush, and "adaptive silence".
Chapter 9. Objections and evidence bearing on the theory of concealing coloration.
In this chapter Cott discusses various objections to the adaptive (evolutionary) nature of camouflage, and provides evidence to dismiss them. Some are "based upon such obvious fallacies that they hardly deserve serious consideration."
Chapter 10. The effectiveness of concealing coloration.
Cott describes simple experiments such as that fish that have changed colour to match a pale background survived better (64% to 42%) on such a background than fish which had not. He also quotes some anecdotal observations on wild animals with similar but not quantified results.
Part II: Advertisement
The methods by which conspicuousness is attained in nature
Chapter 1. The appearance and behaviour of aposematic animals.
Animals that are genuinely distasteful (aposematic) boldly advertise themselves in black, white, red, and yellow. They are often "sluggish", not running from predators; gregarious; and diurnal, since warning displays only work if they can be seen "by potential enemies".
Chapter 2. Warning displays.
Aposematic animals often have (honest) threat displays; edible prey sometimes have (bluffing) startle displays. For example the frilled lizard, Chlamydosaurus kingii, is illustrated in a drawing by Cott, with its tail raised over the body, stretched up on all four legs, mouth wide open, and frills out both sides of the head, making it a startling sight.
Chapter 3. Adventitious warning coloration.
Some marine animals select aposematic materials as coverings, not only as camouflage. Some birds nest near wasps' nests.
Warning coloration in relation to prey
Chapter 4. The nature and function of warning coloration, as illustrated by the mammalia.
Prey like porcupines have warning colours, make noise, and attack predators (even leopards).
Chapter 5. The Protective Attributes Of Aposematic Animals In General.
Evidence is given that conspicuous animals such as caterpillars really are distasteful. Animals with actual poisons are discussed, and how these are secreted, used in bites and stings, or kept to make the animal bitter tasting.
Chapter 6. The relation between warning colours and distasteful attributes.
Various kinds of evidence are presented for aposematism.
Chapter 7. The effectiveness of protective attributes associated with warning colours.
Experimental evidence is presented that insects with warning colours are rejected by predators.
Warning coloration in reference to predatory enemies
Chapter 8. Experimental evidence that vertebrate enemies learn by experience.
Experiments by Cott show that toads learn to avoid eating stinging bees.
Chapter 9. Evidence of selective feeding by vertebrate enemies in a state of nature.
Evidence from wild birds and toads demonstrates preferences for particular prey.
Part III: Disguise
Special protective and aggressive resemblance
Chapter 1. Special resemblance to particular objects.
Cott describes leaf-like fish, chameleons, and insects, and other mimetic forms of camouflage. A liana-like snake near Para (a haunt of Henry Walter Bates in Naturalist on the River Amazons) 160 times as long as it was thick is called "a revelation in the art of aggressive resemblance".
Chapter 2. Adaptive behaviour in relation to special cryptic resemblance.
Animals keep still, sway in the wind, or play dead to assist their camouflage. Poulton's examples of twig-like Geometridae caterpillars are praised. There are fine photographs of leaf insects, and Cott's admired drawing of a poor-me-one or potoo, Nyctibius griseus, sitting on its nest mimicking a broken branch. Cott explains, in a section on "Special resemblances in relation to the attitude of rest"
Chapter 3. Adventitious Concealing Coloration.
Cott begins by citing Shakespeare's Macbeth with "until/ Great Birnamwood to the Dunsinane hill/ Shall come against him" to introduce his chapter on the use of materials as camouflage. Animals from crabs to caterpillars are described.
Conspicuous localized characters
Chapter 4. Deflective marks.
Cott describes markings that help to deflect attack, such as the eyespots of butterfly wings and the twitching cast-off tails of lizards, both acknowledged to Poulton, as well as the distraction displays of birds such as the partridge mentioned by Gilbert White in his Natural History and Antiquities of Selborne.
Chapter 5. Directive marks.
A selection of lures and deceptive markings are described. A large drawing depicts the deimatic warning display of a mantis, Pseudocreobotra wahlbergi with its spined forelegs raised and large spiral eyespots on its spread wings forming an image "suggestive of a formidable foe". Other drawings depict the eyespots of fish such as Chaetodon capistratus, the four-eye butterfly fish, which are "usually towards the tail end" and tending to direct attack away from the head.
Alluring and mimetic resemblances
Chapter 6. Alluring coloration.
The bird-dropping spider Ornithoscatoides decipiens, the flower mantis Hymenopus bicornis and other camouflaged hunters are described.
Chapter 7. Mimicry: the attributes of mimics.
Cott follows Poulton in treating mimicry as basically the same as camouflage or "adaptive resemblance". Batesian mimicry and Mullerian mimicry are compared. The behaviour of "Esquimaux seal-hunters" and First World War Q-ships are mentioned.
Chapter 8. Breeding parasitism and mimicry in cuckoos.
The mimicry and behaviour of the European cuckoo, Cuculus canorus is analysed.
Conclusion
The final chapter confirms that "The force of the facts and arguments used in this work is cumulative in effect." Many small steps of reasoning combine to show that "adaptive coloration... has been... one of the main achievements of organic evolution." The book ends by comparing human artefacts and "natural adaptations", both of which can have goals (recall the publication date of 1940, early in the Second World War) including "the frustration of a predatory animal or of an aggressive Power".
Reception
Foreword
Julian S. Huxley wrote a foreword (labelled 'Introduction') which defends the Darwinian concept of adaptation, especially of colour (in animals) and within that frame of mimicry. He makes it clear that "in these last thirty years" (that is, from about 1910 to 1940) he believed that "experimental biologists" professed, even if they did not actually hold, "a radical scepticism on the subject of adaptations", in other words about whether natural selection really could have created the enormous diversity of pattern and colour seen in nature. Huxley quoted the now long-forgotten Aaron Franklin Shull's 1936 Evolution which stated "These special forms [sexual selection, warning colours, mimicry and signalling] of the selection idea... seem destined to be dropped, or at least relegated to very minor places in the Evolution discussion.", and more sharply that "aggressive and alluring resemblance" (Huxley's words) "must probably be set down as products of fancy belonging to uncritical times." Huxley's reply is simply
With objections dismissed, Huxley remarks that "Dr. Cott is a true follower of Darwin in driving his conclusions home by sheer weight of example," observing that "Faced with his long lists of demonstrative cases, the reader is tempted to wonder why adaptive theories of coloration have been singled out for attack by anti-selectionists." Huxley also noted Cott's "constant cross-reference to human affairs", and that it was good to know that Cott was applying his principles "to the practice of camouflage in war".
Huxley concluded his introduction by describing Adaptive Coloration as "in many respects the last word on the subject", upholding the great tradition of "scientific natural history".
Contemporary reviews (circa 1940)
Reviewers had little to compare Adaptive Coloration with. The English zoologist Edward Bagnall Poulton, a Darwinian, had written a 360-page book, The Colours of Animals, fifty years earlier in 1890, and he was able, at age 84, to review Cott's work in Nature on its appearance in 1940, beginning with the words
The ichthyologist Carl Leavitt Hubbs, reviewing the book for American Naturalist in 1942, began
Hubbs notes that Cott is seeming concerned about the scarcity of experimental data for the survival value of camouflage, and accordingly relies on Sumner and Isely's "clear-cut results", but at once continues that Cott relies on "the general lore of natural history". Hubbs also remarks on the "resurgence to Darwinian views", referring to the scepticism about the power of natural selection among both geneticists of the time and to the Lamarckist views of Trofim Lysenko.
Hubbs observes that Cott is both an artist and a naturalist as well as a scientist: "In section after section, rivaling one another in fascination, this master of art and of natural history unfolds the biological significance of adaptive coloration in animals." And Cott's emphasis on disruptive patterning and (following Thayer) countershading clearly affected the reviewer: "Particularly impressive is the author's treatment of "coincident disruptive coloration", in which a ruptive mark crosses structural boundaries, so as to obliterate visually such ordinarily conspicuous parts as the eye and the limbs. Concealment of an animal's ordinarily telltale shadow is also stressed". Hubbs's review ends "This book is the work of an artist, and it is a work of art. Every biologist with an interest in any phase of natural history or evolution should keep it at hand."
"W.L.S.", reviewing Cott in The Geographical Journal in 1940, begins with "In this large and well-illustrated volume the author discusses at length reason or reasons for the various colour patterns found in the animal kingdom." The reviewer goes on "He has presented us with a vast number of facts and observations which are somewhat difficult to analyse." However "W.L.S." admits that disruptive coloration "is discussed at considerable length by Mr. Cott and many remarkable instances of it are considered in detail". The review ends by mentioning that while biologists (of the 1930s) usually "reject the influence of Natural Selection in evolution, the facts of adaptive coloration as given in Mr. Cott's work are a strong argument in its favour, and must be given due weight. This is what Mr. Cott claims to have accomplished in a volume which will certainly take its place as a most valuable contribution to zoological literature."
Looking back (after 2000)
Peter Forbes, in his book Dazzled and Deceived, wrote that
Over 60 years after its publication, Adaptive Coloration in Animals remains a core reference on the subject. Sören Nylin and colleagues observe in a 2001 paper that
As a natural history narrative on what has become an intensely researched experimental subject, Adaptive Coloration could be thought obsolete, but instead, Peter Forbes observes "But Cott's book is still valuable today for its enormous range, for its passionate exposition of the theories of mimicry and camouflage". This width of coverage and continuing relevance can be seen in the introduction to Sami Merilaita and Johan Lind's 2005 paper on camouflage, Background-Matching and Disruptive Coloration, and the Evolution of Cryptic Coloration, which cites Adaptive Coloration no fewer than eight times, quoting his terms "cryptic coloration or camouflage", "concealing coloration", "background matching (also called cryptic resemblance)", "disruptive coloration", resemblance to visual background, and the difficulty a predator has to detect a prey visually.
Steven Vogel, in a review of Peter Forbes's book Dazzled and Deceived (2009), echoes Julian Huxley's words of seventy years before (in his 'Introduction') by writing
{{blockquote|The zoologist Hugh Cott had the final word in Adaptive Coloration in Animals (1940), a definitive synthesis of everything known about camouflage and mimicry in nature. Cott ruffled fewer feathers [than Trofim Lysenko or Vladimir Nabokov], and his well-organized and unfanatic ideas proved militarily effective, even under the scrutiny of improved techniques for target detection. Thayer’s principles reemerged in more temperate and rational terms, and camouflage schemes based on them survived both photometric analyses and enemy encounters. Biomimetic camouflage took its place as yet another technique in a sophisticated armamentarium of visual deceptions.<ref>Vogel, Steven. The Deceptional Life. American Scientist. On the Bookshelf. September–October 2010. Volume 98, Number 5. Page: 436 </ref>}}
Camouflage researcher Roy Behrens cites and discusses Adaptive Coloration frequently in his writings. For example, in his Camoupedia blog, related to the book of the same name, he writes of Cott's drawings of the hind limbs of the Common frog: "Reproduced above is one of my favorite drawings from what is one of my favorite books." He continues "What makes these drawings (and the book itself) even more interesting is that Cott (1900-1987) was not just a zoologist—he was a highly skilled scientific illustrator (these are his own pen-and-ink drawings), a wildlife photographer, and a prominent British camoufleur in World War II." Still in 2011, Behrens can write of Cott's way of thinking, citing his words as models of clear and accurate explanation of the mechanisms of camouflage: "As he so aptly explained it, disruptive patterns work 'by the optical destruction of what is present', while continuous patterns work 'by the optical construction of what is not present.'"
Publication historyAdaptive Coloration in Animals has been published as follows:
1940, Methuen, Frome and London (printed by Butler and Tanner). Foreword by Julian Huxley
1940, Oxford University Press, New York
1941, Oxford University Press, New York
1957, Methuen, London (reprinted with minor corrections)
1966, Methuen, London (reprinted with minor corrections)
See also
Concealing-Coloration in the Animal Kingdom (G. H. Thayer, 1909)
The Colours of Animals (E. B. Poulton, 1890)
Animal Coloration'' (F. E. Beddard, 1892)
References
Primary
Secondary
Bibliography
External links
Ohio State University: The Camouflage Project: Hugh Cott
Smithsonian: A Painter of Angels Became the Father of Camouflage by Richard Meryman, April 1999 (Article contrasting Thayer and Cott)
Zoology books
Animal coat colors
Mimicry
1940 non-fiction books
Natural history books
Camouflage books | Adaptive Coloration in Animals | [
"Biology"
] | 4,766 | [
"Mimicry",
"Biological defense mechanisms"
] |
37,661,448 | https://en.wikipedia.org/wiki/Katz%E2%80%93Lang%20finiteness%20theorem | In number theory, the Katz–Lang finiteness theorem, proved by , states that if X is a smooth geometrically connected scheme of finite type over a field K that is finitely generated over the prime field, and Ker(X/K) is the kernel of the maps between their abelianized fundamental groups, then Ker(X/K) is finite if K has characteristic 0, and the part of the kernel coprime to p is finite if K has characteristic p > 0.
References
Theorems in number theory | Katz–Lang finiteness theorem | [
"Mathematics"
] | 110 | [
"Number theory stubs",
"Theorems in number theory",
"Mathematical problems",
"Mathematical theorems",
"Number theory"
] |
29,618,699 | https://en.wikipedia.org/wiki/Shearer%27s%20inequality | Shearer's inequality or also Shearer's lemma, in mathematics, is an inequality in information theory relating the entropy of a set of variables to the entropies of a collection of subsets. It is named for mathematician James B. Shearer.
Concretely, it states that if X1, ..., Xd are random variables and S1, ..., Sn are subsets of {1, 2, ..., d} such that every integer between 1 and d lies in at least r of these subsets, then
where is entropy and is the Cartesian product of random variables with indices j in .
Combinatorial version
Let be a family of subsets of [n] (possibly with repeats) with each included in at least members of . Let be another set of subsets of . Then
where the set of possible intersections of elements of with .
See also
Lovász local lemma
References
Information theory
Inequalities | Shearer's inequality | [
"Mathematics",
"Technology",
"Engineering"
] | 201 | [
"Mathematical theorems",
"Telecommunications engineering",
"Applied mathematics",
"Binary relations",
"Computer science",
"Information theory",
"Mathematical relations",
"Inequalities (mathematics)",
"Mathematical problems"
] |
29,621,014 | https://en.wikipedia.org/wiki/R-407C | R-407C is a mixture of hydrofluorocarbons used as a refrigerant. It is a zeotropic blend of difluoromethane (R-32), pentafluoroethane (R-125), and 1,1,1,2-tetrafluoroethane (R-134a). Difluoromethane serves to provide the heat capacity, pentafluoroethane decreases flammability, tetrafluoroethane reduces pressure. R-407C cylinders are colored burnt orange.
This refrigerant is intended as a replacement for R-22. R-22 production will be phased out by 2020 as per the Montreal Protocol.
Key Characteristics:
Colourless and Odourless: R407C is visually clear and does not emit a noticeable smell.
Non-Flammable: Safe to use under normal operating conditions due to its non-flammable nature.
Moderate GWP: While R407C has a lower GWP than R22, it is important to manage its use responsibly to minimize environmental impact.
Applications
R407C is employed in a variety of cooling and refrigeration contexts:
Residential Air Conditioning: Ideal for cooling home environments efficiently.
Commercial Air Conditioning: Suited for larger commercial spaces, including offices and retail locations.
Industrial Refrigeration: Used in industrial processes requiring precise temperature control.
Transport Refrigeration: Essential for maintaining temperatures in refrigerated transport systems.
Features & Benefits
R407C offers several notable benefits:
Replacement for R22: Often used as a direct substitute for R22 in many existing systems, facilitating upgrades with minimal system modifications.
High Efficiency: Provides effective cooling with good energy efficiency.
Environmentally Friendly: Does not deplete the ozone layer, making it a more sustainable option than its predecessors.
Equipment Compatibility: Generally compatible with R22 system components, though some adjustments may be necessary for optimal performance.
Environmental and Safety Considerations
While R407C presents environmental advantages over older refrigerants, it still requires careful handling:
Global Warming Potential: Despite its lower GWP compared to R22, R407C contributes to global warming if released. Effective leak management and system maintenance are essential.
Handling Precautions: R407C should be managed with appropriate safety measures to avoid high-pressure exposure. Personal protective equipment (PPE) and adherence to safety protocols are recommended.
Safety Precautions
For safe handling of R407C, consider the following precautions:
Protective Gear: Wear safety goggles, gloves, and appropriate clothing to prevent exposure.
Ventilation: Ensure proper ventilation in areas where R407C is used to avoid vapor accumulation.
System Maintenance: Regularly check and maintain systems to prevent leaks and ensure efficiency.
Emergency Readiness: Be prepared with emergency procedures in case of accidental exposure or release.
Physical properties
References
Refrigerants
Greenhouse gases | R-407C | [
"Chemistry",
"Environmental_science"
] | 603 | [
"Greenhouse gases",
"Environmental chemistry"
] |
29,623,466 | https://en.wikipedia.org/wiki/Microwave%20digestion | Microwave digestion is a chemical technique used to decompose sample material into a solution suitable for quantitative elemental analysis. It is commonly used to prepare samples for analysis using inductively coupled plasma mass spectrometry (ICP-MS), atomic absorption spectroscopy, and atomic emission spectroscopy (including ICP-AES).
To perform the digestion, sample material is combined with a concentrated strong acid or a mixture thereof, most commonly using nitric acid, hydrochloric acid and/or hydrofluoric acid, in a closed PTFE vessel. The vessel and its contents are then exposed to microwave irradiation, raising the pressure and temperature of the solution mixture. The elevated pressures and temperatures within a low pH sample medium increase both the speed of thermal decomposition of the sample and the solubility of elements in solution. Organic compounds are decomposed into gaseous form, effectively removing them from solution. Once these elements are in solution, it is possible to quantify elemental concentrations within samples.
Microwaves can be programmed to reach specific temperatures or ramp up to a given temperature at a specified rate. The temperature in the interior of the vessel is monitored by an infrared external sensor or by a optic fiber probe, and the microwave power is regulated to maintain the temperature defined by the active program. The vessel solution must contain at least one solvent that absorbs microwave radiation, usually water. The specific blend of acids (or other reagents) and the temperatures vary depending upon the type of sample being digested. Often a standardized protocol for digestion is followed, such as an Environmental Protection Agency Method.
Comparison between microwave digestion and other sample preparation methods
Before microwave digestion technology was developed, samples were digested using less convenient methods, such as heating vessels in an oven, typically for at least 24 hours. The use of microwave energy allows for fast sample heating, reducing digestion time to as little as one hour.
Another common means to decompose samples for elemental analysis is dry-ashing, in which samples are incinerated in a muffle furnace. The resultant ash is then dissolved for analysis, usually into dilute nitric acid. While this method is simple, inexpensive and does not require concentrated acids, it cannot be used for volatile elements such as mercury and can increase the likelihood of background contamination. The incineration will not convert all elements to soluble salts, necessitating an additional digestion step.
Quality control in microwave digestion
In microwave digestion, 100% analyte recovery cannot be assumed. To account for this, scientists perform tests such as fortification recovery, in which a spike (a known amount of the target analyte) is added to test samples. These spiked samples are then analyzed to determine whether the expected increase in analyte concentration occurs.
Contamination from improperly cleaned digestion vessels is also a possibility. As such, in any microwave digestion, blank samples need to be digested to determine if there is background contamination.
References
Footnotes
Bibliography
Analytical chemistry | Microwave digestion | [
"Chemistry"
] | 608 | [
"nan"
] |
29,630,460 | https://en.wikipedia.org/wiki/OD600 | OD600 (Also written as O.D. 600, D600, o.d. 600, OD600) is an abbreviation indicating the optical density of a sample measured at a wavelength of 600 nm in 1 cm light path (unless otherwise stated). It is a commonly used in microbiology for estimating the concentration of bacteria or other cells in a liquid as the 600 nm wavelength does little to damage or hinder their growth. OD600 is a type of turbidity measurement. Since optical density in case of OD600 measurements results from light scattering by particles (cells) rather than absorption, size and shape as well as dead cells and debris of a cell may add to light dissipating. Distinctive cell types that are at densities of the same level (eg. cell/mL), may, therefore, show varying values OD600, when estimated on a similar instrument.
For turbid samples such as cell cultures, the major contributor for the optical density measured is light scattering and not the result of molecular absorption following the Beer-Lambert Law. The measurements are therefore depending on the optical setup of the spectrophotometer (distance between the cell holder and instrument exit slit, monochromator optics, slit geometry, etc.), different instrument types will most likely tend to give different OD600 readings for the same turbid sample.
Measuring the change of the OD600 as a function of time (e.g. measuring of the growth curve) may indicate the growth phase of cultured cell population, i.e. whether it is in the lag phase, log phase, or stationary phase. For protein expression and purification in bacteria it is recommended that protein induction and cell harvesting should be done at specific OD600 (usually at the end of the log phase, OD600 = 0.4).
OD600 is preferable to UV spectroscopy when measuring the growth over time of a cell population because at this wavelength, the cells will not be killed as they would under too much UV radiation. UV radiation has also been shown to cause small to medium-sized mutations in bacteria, potentially altering or destroying genes of interest.
References
Absorption spectroscopy | OD600 | [
"Physics",
"Chemistry"
] | 462 | [
"Spectroscopy",
"Spectrum (physical sciences)",
"Absorption spectroscopy"
] |
29,631,337 | https://en.wikipedia.org/wiki/Benzyltrimethylammonium%20hydroxide | Benzyltrimethylammonium hydroxide, also known as Triton B or trimethylbenzylammonium hydroxide, is a quaternary ammonium salt that functions as an organic base. It is usually handled as a solution in water or methanol. The compound is colourless, although the solutions often appear yellowish. Commercial samples often have a distinctive fish-like odour, presumably due to the presence of trimethylamine via hydrolysis.
Uses
Together with the benzyltriethylammonium salt, benzyltrimethylammonium hydroxide is a popular phase-transfer catalyst.
It is used in aldol condensation reactions and base-catalyzed dehydration reactions. It is also used as a base in Ando's Z-selective variant of Horner-Wadsworth-Emmons Olefination reactions.
Relative to tetramethylammonium hydroxide, benzyltriethylammonium hydroxide is more labile. In 6M NaOH at 160 °C their half-lives are 61.9 and 4 h, respectively.
References
See also
http://www.chemcas.com/AnalyticalDetail.asp?pidx=1&id=19594&cas=100-85-6&page=490
Hydroxides
Quaternary ammonium compounds
Reagents for organic chemistry
Benzyl compounds | Benzyltrimethylammonium hydroxide | [
"Chemistry"
] | 290 | [
"Bases (chemistry)",
"Hydroxides",
"Reagents for organic chemistry"
] |
29,631,617 | https://en.wikipedia.org/wiki/Antimatter%20comet | Antimatter comets and antimatter meteoroids are hypothetical comets and meteoroids composed solely of antimatter instead of ordinary matter. Although never actually observed, and unlikely to exist anywhere within the Milky Way, they have been hypothesized to exist, and their existence, on the presumption that hypothesis is correct, has been put forward as one possible explanation for various observed natural phenomena over the years.
Hypothesized existence
The hypothesis of comets made of antimatter can be traced back to the 1940s, when physicist Vladimir Rojansky proposed, in his paper "The Hypothesis of the Existence of Contraterrene Matter", the possibility that some comets and meteoroids could be made from "contraterrene" matter (i.e. antimatter). Such objects, Rojanski stated, would (if they existed at all) have their origins outside the Solar System. He hypothesized that if there were an antimatter object in orbit in the Solar System, it would exhibit the behavior of comets observed in the 1940s: As its atoms annihilated with "terrene" matter from other bodies and solar wind, it would generate volatile compounds and undergo a change of composition to elements with lower atomic masses. From this basis he propounded the hypothesis that some objects that had been identified as comets may, in fact, be antimatter objects, suggesting, based upon calculations using the Stefan–Boltzmann law, that it would be possible to determine the existence of such objects within the Solar System by observing their temperatures. An antimatter body subjected to normal levels of meteoric bombardment (per 1940s figures), and absorbing half of the energy created by the annihilation of normal matter and antimatter, would have a temperature of for bombardment figures calculated by Wylie or for calculations by Nininger. In the 1970s, when comet Kohoutek was observed, Rojanski again suggested hypothesis of antimatter comets in a letter in Physical Review Letters, and suggested that gamma-ray observations be made of the comet to test this hypothesis.
Rojansky's original 1940 hypothesis was that perhaps the only bodies within the Solar System that could be antimatter were comets and meteoroids, all others being almost certainly normal matter. Experimental evidence gathered since then has not only borne out this restriction but has made the existence of actual antimatter comets and meteoroids themselves seem ever more unlikely. Gary Steigman, assistant professor of Astronomy at Yale University, observed in 1976 that space probes had proven — by the fact that they were not annihilated upon impact — that bodies such as Mars, Venus, and the Moon were not antimatter. He also noted that had any of the planets or similar bodies been antimatter, their interaction with the terrene solar wind and the sheer strength of the gamma ray emissions that would have resulted would have made them readily noticeable long since. He noted that not even antimatter cosmic rays had been found, with all of the nuclei found in studies having been uniformly terrene, the experimental data in several studies made from 1961 onwards by various people excluding the presence of a fractional antimatter composition of cosmic rays any larger than 10−4 of the total. Further, the uniformly terrene nature of the cosmic ray flux indicates that nowhere in the Milky Way are there any sources of heavier antimatter elements (such as carbon), since (although it is not proven) it is a likely assumption that they represent the overall composition of the entire galaxy. They are representative of the galaxy as a whole — goes the logic — and since they do contain terrene carbon and other atoms, but have not been observed to contain any antimatter atoms, therefore there is no reasonable source for extrasolar antimatter comets, meteoroids, or any other large scale heavy element objects to originate from, within this galaxy.
Martin Beech from the University of Western Ontario (London, Ontario, Canada) referred to the various hypotheses and experimental results that support non-existence of antimatter in the Universe. He argued that any antimatter comets and meteors that exist must be (at least) extrasolar in origin because the nebular hypothesis for the formation of the Solar System precludes their being solar. Any antimatter in a pre-formation nebula or planetary accretion disc has a comparatively short lifetime, in astronomical terms, before annihilation with the terrene matter that it is mixed with. This lifetime is measured in the hundreds of years, and so any solar antimatter present at the time that the system was formed will have long since been annihilated. Any antimatter comets and meteors must therefore come from another solar system. Furthermore, not only must antimatter meteors be extrasolar in origin, they must have been recently (i.e. within the past 104 ~ 105 years) captured by the Solar System. Most meteoroids are broken down to sizes of 10−5 g within that timeframe, because of meteoroid-upon-meteoroid collisions. Thus any antimatter meteor must be either extrasolar in origin itself, or broken off from an antimatter comet that is extrasolar in origin. The former are unlikely to exist from observational evidence. Any extrasolar meteoroid would have a hyperbolic orbit, but less than 1% of the observed meteoroids have such, and the process of perturbation of ordinary (terrene) solar objects, by planetary encounters, into hyperbolic trajectories accounts for all of those. Beech concluded that a continued null result, however, does not constitute a proof ('Absence of evidence is not evidence of absence', M. Rees) and a single positive detection negates the arguments presented.
Hypothesized explanations for observed phenomena
Tektites
In 1947, Mohammad Abdur Rahman Khan, professor at Osmania University and research associate at the Institute of Meteoretics in the University of New Mexico, put forward the hypothesis that antimatter comets or meteoroids were responsible for tektites . However, this explanation, out of the many proposed explanations for tektites, is considered to be one of the more improbable.
Tunguska event of 1908
By the 1950s, speculating about antimatter comets and meteoroids was a commonplace exercise for astrophysicists. One such, Philip J. Wyatt of Florida State University, suggested that the Tunguska event may have been a meteor made of antimatter . Willard Libby and Clyde Cowan took Wyatt's idea further , having studied worldwide levels of carbon-14 in tree rings and noticing unusually high levels for the year 1909. However, even in 1958 the theoretical flaws in the hypothesis were observed, aside from the evidence that was coming in at the same time from the first gamma ray measurement satellites. For one, the hypothesis did not explain how an antimatter meteor could have managed to survive that low into the Earth's atmosphere, without being annihilated as soon as it encountered terrene matter at the upper levels.
Ball lightning
In 1971, fragments of antimatter comets or meteoroids were hypothesized, by David E. T. F. Ashby of Culham Laboratory and Colin Whitehead of the U.K. Atomic Energy Research Establishment, as a possible cause for ball lightning . They monitored the sky with gamma-ray detection apparatus, and reported unusually high numbers at 511 keV (kilo-electron volts) which is the characteristic gamma ray frequency of a collision between an electron and a positron. There were natural explanations for such readings. In particular positrons can be produced indirectly by the action of a thunderstorm, as it creates the unstable isotopes nitrogen-13 and oxygen-15. However, Ashby and Whitehead noted that there were no thunderstorms present at the times that the gamma-ray readings were observed. They instead presented the hypothesis of antimatter meteors as an interesting one that did explain all of what their observations had recorded, and suggested that it merited further investigation.
Gamma-ray bursts
Antimatter comets thought to exist in the Oort cloud were in the 1990s hypothesized as one possible explanation for gamma-ray bursts. These bursts can be explained by the annihilation of matter and antimatter microcomets. The explosion would create powerful gamma ray bursts and accelerate matter to near light speeds. These antimatter microcomets are thought to reside at distances of more than 1000 AU. Calculations have shown that comets of around 1 km in radius would shrink by 1 m if they passed the Sun with a perihelion of 1 AU. Microcomets, due to the stresses of solar heating, shatter and burn up much more quickly because the forces are more concentrated within their small masses. Antimatter microcomets would burn up even more rapidly because the annihilation of solar wind with the surface of the microcomet would produce additional heat. As more gamma-ray bursts were detected in subsequent years, this theory failed to explain the observed distribution of gamma-ray bursts about host galaxies and detections of X-ray lines associated with gamma-ray bursts. The discovery of a supernova associated with a gamma-ray burst in 2002 provided compelling evidence that massive stars are the origin of gamma-ray bursts. Since 2002, more supernovae have been observed to be associated with gamma-ray bursts, and massive stars as the origin of gamma-ray bursts has been firmly established.
Footnotes
References
Bibliography
Further reading
Original publications of the various hypotheses
Other
English translation:
Antimatter
Comets
Meteoroids
Hypothetical astronomical objects
Tunguska event | Antimatter comet | [
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"Astronomy"
] | 1,976 | [
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"Antimatter",
"Unsolved problems in physics",
"Astronomical myths",
"Hypothetical astronomical objects",
"Tunguska event",
"Astronomical objects",
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46,464,450 | https://en.wikipedia.org/wiki/Averaged%20Lagrangian | In continuum mechanics, Whitham's averaged Lagrangian method – or in short Whitham's method – is used to study the Lagrangian dynamics of slowly-varying wave trains in an inhomogeneous (moving) medium.
The method is applicable to both linear and non-linear systems. As a direct consequence of the averaging used in the method, wave action is a conserved property of the wave motion. In contrast, the wave energy is not necessarily conserved, due to the exchange of energy with the mean motion. However the total energy, the sum of the energies in the wave motion and the mean motion, will be conserved for a time-invariant Lagrangian. Further, the averaged Lagrangian has a strong relation to the dispersion relation of the system.
The method is due to Gerald Whitham, who developed it in the 1960s. It is for instance used in the modelling of surface gravity waves on fluid interfaces, and in plasma physics.
Resulting equations for pure wave motion
In case a Lagrangian formulation of a continuum mechanics system is available, the averaged Lagrangian methodology can be used to find approximations for the average dynamics of wave motion – and (eventually) for the interaction between the wave motion and the mean motion – assuming the envelope dynamics of the carrier waves is slowly varying. Phase averaging of the Lagrangian results in an averaged Lagrangian, which is always independent of the wave phase itself (but depends on slowly varying wave quantities like wave amplitude, frequency and wavenumber). By Noether's theorem, variation of the averaged Lagrangian with respect to the invariant wave phase then gives rise to a conservation law:
This equation states the conservation of wave action – a generalization of the concept of an adiabatic invariant to continuum mechanics – with
being the wave action and wave action flux respectively. Further and denote space and time respectively, while is the gradient operator. The angular frequency and wavenumber are defined as
and both are assumed to be slowly varying. Due to this definition, and have to satisfy the consistency relations:
The first consistency equation is known as the conservation of wave crests, and the second states that the wavenumber field is irrotational (i.e. has zero curl).
Method
The averaged Lagrangian approach applies to wave motion – possibly superposed on a mean motion – that can be described in a Lagrangian formulation. Using an ansatz on the form of the wave part of the motion, the Lagrangian is phase averaged. Since the Lagrangian is associated with the kinetic energy and potential energy of the motion, the oscillations contribute to the Lagrangian, although the mean value of the wave's oscillatory excursion is zero (or very small).
The resulting averaged Lagrangian contains wave characteristics like the wavenumber, angular frequency and amplitude (or equivalently the wave's energy density or wave action). But the wave phase itself is absent due to the phase averaging. Consequently, through Noether's theorem, there is a conservation law called the conservation of wave action.
Originally the averaged Lagrangian method was developed by Whitham for slowly-varying dispersive wave trains. Several extensions have been made, e.g. to interacting wave components, Hamiltonian mechanics, higher-order modulational effects, dissipation effects.
Variational formulation
The averaged Lagrangian method requires the existence of a Lagrangian describing the wave motion. For instance for a field , described by a Lagrangian density the principle of stationary action is:
with the gradient operator and the time derivative operator. This action principle results in the Euler–Lagrange equation:
which is the second-order partial differential equation describing the dynamics of Higher-order partial differential equations require the inclusion of higher than first-order derivatives in the Lagrangian.
Example
For example, consider a non-dimensional and non-linear Klein–Gordon equation in one space dimension :
This Euler–Lagrange equation emerges from the Lagrangian density:
The small-amplitude approximation for the Sine–Gordon equation corresponds with the value For the system is linear and the classical one-dimensional Klein–Gordon equation is obtained.
Slowly-varying waves
Slowly-varying linear waves
Whitham developed several approaches to obtain an averaged Lagrangian method. The simplest one is for slowly-varying linear wavetrains, which method will be applied here.
The slowly-varying wavetrain –without mean motion– in a linear dispersive system is described as:
with and
where is the real-valued wave phase, denotes the absolute value of the complex-valued amplitude while is its argument and denotes its real part. The real-valued amplitude and phase shift are denoted by and respectively.
Now, by definition, the angular frequency and wavenumber vector are expressed as the time derivative and gradient of the wave phase as:
and
As a consequence, and have to satisfy the consistency relations:
and
These two consistency relations denote the "conservation of wave crests", and the irrotationality of the wavenumber field.
Because of the assumption of slow variations in the wave train – as well as in a possible inhomogeneous medium and mean motion – the quantities and all vary slowly in space and time – but the wave phase itself does not vary slowly. Consequently, derivatives of and are neglected in the determination of the derivatives of for use in the averaged Lagrangian:
and
Next these assumptions on and its derivatives are applied to the Lagrangian density
Slowly-varying non-linear waves
Several approaches to slowly-varying non-linear wavetrains are possible. One is by the use of Stokes expansions, used by Whitham to analyse slowly-varying Stokes waves. A Stokes expansion of the field can be written as:
where the amplitudes etc. are slowly varying, as are the phases etc. As for the linear wave case, in lowest order (as far as modulational effects are concerned) derivatives of amplitudes and phases are neglected, except for derivatives and of the fast phase :
and
These approximations are to be applied in the Lagrangian density , and its phase average
Averaged Lagrangian for slowly-varying waves
For pure wave motion the Lagrangian is expressed in terms of the field and its derivatives. In the averaged Lagrangian method, the above-given assumptions on the field –and its derivatives– are applied to calculate the Lagrangian. The Lagrangian is thereafter averaged over the wave phase :
As a last step, this averaging result can be expressed as the averaged Lagrangian density – which is a function of the slowly varying parameters and and independent of the wave phase itself.
The averaged Lagrangian density is now proposed by Whitham to follow the average variational principle:
From the variations of follow the dynamical equations for the slowly-varying wave properties.
Example
Continuing on the example of the nonlinear Klein–Gordon equation, see equations and , and applying the above approximations for and (for this 1D example) in the Lagrangian density, the result after averaging over is:
where it has been assumed that, in big-O notation, and . Variation of with respect to leads to So the averaged Lagrangian is:
For linear wave motion the averaged Lagrangian is obtained by setting equal to zero.
Set of equations emerging from the averaged Lagrangian
Applying the averaged Lagrangian principle, variation with respect to the wave phase leads to the conservation of wave action:
since and while the wave phase does not appear in the averaged Lagrangian density due to the phase averaging.
Defining the wave action as and the wave action flux as the result is:
The wave action equation is accompanied with the consistency equations for and which are:
and
Variation with respect to the amplitude leads to the dispersion relation
Example
Continuing with the nonlinear Klein–Gordon equation, using the average variational principle on equation , the wave action equation becomes by variation with respect to the wave phase
and the nonlinear dispersion relation follows from variation with respect to the amplitude
So the wave action is and the wave action flux The group velocity is
Mean motion and pseudo-phase
Conservation of wave action
The averaged Lagrangian is obtained by integration of the Lagrangian over the wave phase. As a result, the averaged Lagrangian only contains the derivatives of the wave phase (these derivatives being, by definition, the angular frequency and wavenumber) and does not depend on the wave phase itself. So the solutions will be independent of the choice of the zero level for the wave phase. Consequently – by Noether's theorem – variation of the averaged Lagrangian with respect to the wave phase results in a conservation law:
where
with the wave action and the wave action flux. Further denotes the partial derivative with respect to time, and is the gradient operator. By definition, the group velocity is given by:
Note that in general the energy of the wave motion does not need to be conserved, since there can be an energy exchange with a mean flow. The total energy – the sum of the energies of the wave motion and the mean flow – is conserved (when there is no work by external forces and no energy dissipation).
Conservation of wave action is also found by applying the generalized Lagrangian mean (GLM) method to the equations of the combined flow of waves and mean motion, using Newtonian mechanics instead of a variational approach.
Conservation of energy and momentum
Connection to the dispersion relation
Pure wave motion by linear models always leads to an averaged Lagrangian density of the form:
Consequently, the variation with respect to amplitude: gives
So this turns out to be the dispersion relation for the linear waves, and the averaged Lagrangian for linear waves is always the dispersion function times the amplitude squared.
More generally, for weakly nonlinear and slowly modulated waves propagating in one space dimension and including higher-order dispersion effects – not neglecting the time and space derivatives and of the amplitude when taking derivatives, where is a small modulation parameter – the averaged Lagrangian density is of the form:
with the slow variables and
References
Notes
Publications by Whitham on the method
An overview can be found in the book:
Some publications by Whitham on the method are:
Further reading
Continuum mechanics
Lagrangian mechanics | Averaged Lagrangian | [
"Physics",
"Mathematics"
] | 2,160 | [
"Dynamical systems",
"Lagrangian mechanics",
"Classical mechanics",
"Continuum mechanics"
] |
46,464,527 | https://en.wikipedia.org/wiki/Software%20Guard%20Extensions | Intel Software Guard Extensions (SGX) is a set of instruction codes implementing trusted execution environment that are built into some Intel central processing units (CPUs). They allow user-level and operating system code to define protected private regions of memory, called enclaves. SGX is designed to be useful for implementing secure remote computation, secure web browsing, and digital rights management (DRM). Other applications include concealment of proprietary algorithms and of encryption keys.
SGX involves encryption by the CPU of a portion of memory (the enclave). Data and code originating in the enclave are decrypted on the fly within the CPU, protecting them from being examined or read by other code, including code running at higher privilege levels such as the operating system and any underlying hypervisors. While this can mitigate many kinds of attacks, it does not protect against side-channel attacks.
A pivot by Intel in 2021 resulted in the deprecation of SGX from the 11th and 12th generation Intel Core processors, but development continues on Intel Xeon for cloud and enterprise use.
Details
SGX was first introduced in 2015 with the sixth generation Intel Core microprocessors based on the Skylake microarchitecture.
Support for SGX in the CPU is indicated in CPUID "Structured Extended feature Leaf", EBX bit 02, but its availability to applications requires BIOS/UEFI support and opt-in enabling which is not reflected in CPUID bits. This complicates the feature detection logic for applications.
Emulation of SGX was added to an experimental version of the QEMU system emulator in 2014. In 2015, researchers at the Georgia Institute of Technology released an open-source simulator named "OpenSGX".
One example of SGX used in security was a demo application from wolfSSL using it for cryptography algorithms.
Intel Goldmont Plus (Gemini Lake) microarchitecture also contains support for Intel SGX.
Both in the 11th and 12th generations of Intel Core processors, SGX is listed as "Deprecated" and thereby not supported on "client platform" processors. This removed support of playing Ultra HD Blu-ray discs on officially licensed software, such as PowerDVD.
List of SGX vulnerabilities
Prime+Probe attack
On 27 March 2017 researchers at Austria's Graz University of Technology developed a proof-of-concept that can grab RSA keys from SGX enclaves running on the same system within five minutes by using certain CPU instructions in lieu of a fine-grained timer to exploit cache DRAM side-channels. One countermeasure for this type of attack was presented and published by Daniel Gruss et al. at the USENIX Security Symposium in 2017. Among other published countermeasures, one countermeasure to this type of attack was published on September 28, 2017, a compiler-based tool, DR.SGX, that claims to have superior performance with the elimination of the implementation complexity of other proposed solutions.
Spectre-like attack
The LSDS group at Imperial College London showed a proof of concept that the Spectre speculative execution security vulnerability can be adapted to attack the secure enclave. The Foreshadow attack, disclosed in August 2018, combines speculative execution and buffer overflow to bypass the SGX. A security advisory and mitigation for this attack, also called an L1 Terminal Fault, was originally issued on August 14, 2018 and updated May 11, 2021.
Enclave attack
On 8 February 2019, researchers at Austria's Graz University of Technology published findings which showed that in some cases it is possible to run malicious code from within the enclave itself. The exploit involves scanning through process memory in order to reconstruct a payload, which can then run code on the system. The paper claims that due to the confidential and protected nature of the enclave, it is impossible for antivirus software to detect and remove malware residing within it. Intel issued a statement, stating that this attack was outside the threat model of SGX, that they cannot guarantee that code run by the user comes from trusted sources, and urged consumers to only run trusted code.
MicroScope replay attack
There is a proliferation of side-channel attacks plaguing modern computer architectures. Many of these attacks measure slight, nondeterministic variations in the execution of code, so the attacker needs many measurements (possibly tens of thousands) to learn secrets. However, the MicroScope attack allows a malicious OS to replay code an arbitrary number of times regardless of the program's actual structure, enabling dozens of side-channel attacks. In July 2022, Intel submitted a Linux patch called AEX-Notify to allow the SGX enclave programmer to write a handler for these types of events.
Plundervolt
Security researchers were able to inject timing specific faults into execution within the enclave, resulting in leakage of information. The attack can be executed remotely, but requires
access to the privileged control of the processor's voltage and frequency. A security advisory and mitigation for this attack was originally issued on August 14, 2018 and updated on March 20, 2020.
LVI
Load Value Injection injects data into a program aiming to replace the value loaded from memory which is then used for a short time before the mistake is spotted and rolled back, during which LVI controls data and control flow. A security advisory and mitigation for this attack was originally issued on March 10, 2020 and updated on May 11, 2021.
SGAxe
SGAxe, an SGX vulnerability published in 2020, extends a speculative execution attack on cache, leaking content of the enclave. This allows an attacker to access private CPU keys used for remote attestation. In other words, a threat actor can bypass Intel's countermeasures to breach SGX enclaves' confidentiality. The SGAxe attack is carried out by extracting attestation keys from SGX's private quoting enclave that are signed by Intel. The attacker can then masquerade as legitimate Intel machines by signing arbitrary SGX attestation quotes. A security advisory and mitigation for this attack, also called a Processor Data Leakage or Cache Eviction, was originally issued January 27, 2020 and updated May 11, 2021.
ÆPIC leak
In 2022, security researchers discovered a vulnerability in the Advanced Programmable Interrupt Controller (APIC) that allows for an attacker with root/admin privileges to gain access to encryption keys via the APIC by inspecting data transfers from L1 and L2 cache. This vulnerability is the first architectural attack discovered on x86 CPUs. This differs from Spectre and Meltdown which use a noisy side channel. This exploit currently affects Intel Core 10th, 11th and 12th generations, and Xeon Ice Lake microprocessors.
Extraction of the private key
The code signature is generated with a private key that is only in the enclave. The private key is encoded via “fuse” elements on the chip. In the process, bits are burnt through, giving them the binary value 0. This private key cannot be extracted because it is encoded in the hardware. Mark Ermolov, Maxim Goryachy and Dmitry Sklyarov refuted the claim to trustworthiness of the SGX concept at https://github.com/chip-red-pill/glm-ucode#.
SGX malware arguments
There has been a long debate on whether SGX enables creation of superior malware. Oxford University researchers published an article in October 2022 considering attackers' potential advantages and disadvantages by abusing SGX for malware development. Researchers conclude that while there might be temporary zero-day vulnerabilities to abuse in SGX ecosystem, the core principles and design features of Trusted Execution Environments (TEEs) make malware weaker than a malware-in-the-wild, TEEs make no major contributions to malware otherwise.
See also
Intel MPX
Spectre-NG
Speculative execution CPU vulnerabilities
References
External links
Intel Software Guard Extensions (Intel SGX) / ISA Extensions, Intel
Intel Software Guard Extensions (Intel SGX) Programming Reference, Intel, October 2014
IDF 2015 - Tech Chat: A Primer on Intel Software Guard Extensions, Intel (poster)
ISCA 2015 tutorial slides for Intel SGX, Intel, June 2015
McKeen, Frank, et al. (Intel), Innovative Instructions and Software Model for Isolated Execution // Proceedings of the 2nd International Workshop on Hardware and Architectural Support for Security and Privacy. ACM, 2013.
Jackson, Alon, (PhD dissertation). Trust is in the Keys of the Beholder: Extending SGX Autonomy and Anonymity, May 2017.
Joanna Rutkowska, Thoughts on Intel's upcoming Software Guard Extensions (Part 1), August 2013
SGX: the good, the bad and the downright ugly / Shaun Davenport, Richard Ford (Florida Institute of Technology) / Virus Bulletin, 2014-01-07
Victor Costan and Srinivas Devadas, Intel SGX Explained, January 2016.
wolfSSL, October 2016.
The Security of Intel SGX for Key Protection and Data Privacy Applications / Professor Yehuda Lindell (Bar Ilan University & Unbound Tech), January 2018
Intel SGX Technology and the Impact of Processor Side-Channel Attacks, March 2020
How Confidential Computing Delivers A Personalised Shopping Experience, January 2021
Realising the Potential of Data Whilst Preserving Privacy with EyA and Conclave from R3, December 2021
Introduction to Intel Software Guard Extensions, June 2020
Intel
X86 instructions
Cybersecurity engineering | Software Guard Extensions | [
"Technology",
"Engineering"
] | 1,956 | [
"Cybersecurity engineering",
"Computer networks engineering",
"Computer engineering"
] |
43,334,966 | https://en.wikipedia.org/wiki/Valvulotome | A valvulotome is a catheter-based controllable surgical instrument used for cutting or disabling the venous valves. This is needed to enable an in situ bypass in patients with an occluded artery (especially femoral artery), where the saphenous vein is disconnected from the venous system and connected to arteries above and below the occluded segment to allow blood to flow to the lower leg.
Since the leg veins usually contain a number of valves that direct flow towards the heart, they cannot directly be used as graft, but if vein valves are removed the arterial blood can flow via the GSV to the lower leg - this is called an in situ graft procedure, a type of vascular bypass.
The valvulotome itself is a long, flexible catheter with a recessed cutting blade at its end for the destruction of venous valves. The valvulotome is inserted at the distal end of the vein, guided to the proximal end, then withdrawn. It is during withdrawal that the valves are destroyed. The blade is designed to prevent exposure of the vein intima to the sharp cutting surface to avoid damage to the vessel wall. It is often designed resembling a hook, with a blunt outer surface and a sharp inner surface that makes contact with the venous valve as the device is withdrawn, but not during insertion.
References
Medical devices | Valvulotome | [
"Biology"
] | 286 | [
"Medical devices",
"Medical technology"
] |
43,336,040 | https://en.wikipedia.org/wiki/Paul%20Connett | Paul Connett is a prominent water fluoridation critic, executive director of the Binghamton, New York based Fluoride Action Network (FAN), one of the largest organizations opposing water fluoridation worldwide. Critics have stated that The Fluoride Action Network is funded, at least in part, by Joseph Mercola, who has been identified by the Centre for Countering Digital Hate as a leading purveyor of COVID-19 disinformation. FAN executive director Stuart Cooper has stated, "Mercola is among thousands of donors and his money accounts for a single-digit percentage of FAN's contributions".
Connett has been invited by environmental organizations opposing fluoridation to lecture on the subject in fluoridating countries such as Canada, Israel, Australia and New Zealand. Connett has stated, "It’s politics that is interfering with science in this issue...It’s a matter of political will, and you cannot change political will if you don’t get the people. We must involve the people."
Early life
Connett is English, and is a graduate of Cambridge University. He holds a Ph.D. in chemistry from Dartmouth College.
Political activism
Connett became involved in political activism in 1968. He volunteered for Eugene McCarthy's presidential campaign, where he met Allard K. Lowenstein. Lowenstein and Connett founded the American Committee to Keep Biafra Alive, in response to the famine caused by the Biafran War. In 1971, Connett co-founded Operation Omega, a non-violent group taking humanitarian aid into East Pakistan during the Bangladesh Liberation War. Paul's wife Ellen was arrested during one of Omega's trips into East Pakistan and spent two months imprisoned there.
University career
After teaching chemistry and toxicology for 23 years at St. Lawrence University, Canton, NY, he retired from his full professorship. He is currently also the director of the American Environmental Health Studies Project (AEHSP).
Fluoridation campaign
In 2004, Connett published the paper 50 Reasons to Oppose Fluoridation in Medical Veritas, a pseudoscientific journal described by QuackWatch as "fundamentally flawed". In 2010 he coauthored; The Case against Fluoride: How Hazardous Waste Ended Up in Our Drinking Water and the Bad Science and Powerful Politics That Keep It There along with Dr. James Beck and Dr. H. Spedding Micklem.
He also wrote the book in 2013; "The Zero Waste Solution". and assisted the city of Naples in pursuing its zero waste strategy.
References
Year of birth missing (living people)
Living people
Water fluoridation
21st-century American chemists
Alumni of the University of Cambridge
Dartmouth College alumni | Paul Connett | [
"Chemistry"
] | 553 | [
"Water treatment",
"Water fluoridation"
] |
43,341,147 | https://en.wikipedia.org/wiki/Byerlee%27s%20law | In rheology, Byerlee's law, also known as Byerlee's friction law concerns the shear stress (τ) required to slide one rock over another. The rocks have macroscopically flat surfaces, but the surfaces have small asperities that make them "rough." For a given experiment and at normal stresses (σn) below about 2000 bars (200 MPa) the shear stress increases approximately linearly with the normal stress (τ = 0.85 σn, where τ and σn is in units of MPa) and is highly dependent on rock type and the character (roughness) of the surfaces, see Mohr-Coulomb friction law. Byerlee's law states that with increased normal stress the required shear stress continues to increase, but the rate of increase decreases (τ = 50 + 0.6σn), where τ and σn are in units of MPa, and becomes nearly independent of rock type.
The law describes an important property of crustal rock, and can be used to determine when slip along a geological fault takes place.
The law is named after the American geophysicist James Byerlee, who derived it experimentally in 1978.
See also
Archie's law
Birch's law
References
Inline citations
General references
Geophysics
Rheology
Scientific laws | Byerlee's law | [
"Physics",
"Chemistry",
"Mathematics"
] | 274 | [
"Applied and interdisciplinary physics",
"Mathematical objects",
"Scientific laws",
"Equations",
"Geophysics",
"Rheology",
"Fluid dynamics"
] |
28,235,244 | https://en.wikipedia.org/wiki/Bungaroosh |
Bungaroosh (also spelt bungeroosh and other variations) is a composite building material used almost exclusively in the English seaside resort of Brighton, the neighbouring town of Hove and in the coastal Sussex area. The etymology of the word is unknown. Its use dates from the start of the Regency period at the end of the 18th century, and into the 19th when Brighton grew from a fishing village into a large town. Bungaroosh is often found in buildings of that era in the town and in its near neighbours Worthing and Lewes but is little known elsewhere except London. It was a building material first introduced by the Romans and has characteristics of that era. It can incorporate any of a wide variety of substances and materials and is used most often in external walls.
The manufacture of bungaroosh involved placing miscellaneous materials, such as whole or broken bricks, cobblestones, flints (commonly found on the South Downs around Brighton), small pebbles, sand and pieces of wood into hydraulic lime and then by shovelling it between shuttering until it has set. Other structural fittings, such as brick piers or wooden lintels, could then be added if more support was needed. This was particularly common in Brighton where bungaroosh walls were often built behind the stuccoed façades of Regency-style houses. Another technique was to wait for the mixture to set, then render it with a lime-based mixture and paint it. This produced a consistent, regular surface which could be used to build the symmetrical façades required in Georgian architecture, a popular style in Lewes. The material is particularly prevalent in the early 19th-century squares, crescents and terraces of Brighton's seafront, such as Regency Square, Royal Crescent and the Kemp Town estate.
See also
Buildings and architecture of Brighton and Hove
Core-and-veneer
References
Notes
Bibliography
Other resources
(from The Regency Society)
Composite materials
Building materials
Types of wall
Brighton | Bungaroosh | [
"Physics",
"Engineering"
] | 394 | [
"Structural engineering",
"Building engineering",
"Composite materials",
"Construction",
"Types of wall",
"Materials",
"Building materials",
"Matter",
"Architecture"
] |
49,349,799 | https://en.wikipedia.org/wiki/BMJ%20Health%20%26%20Care%20Informatics | BMJ Health & Care Informatics is a peer-reviewed open-access medical journal covering health informatics. It was established in 1992 and is published by the BMJ. The editor-in-chief is Yu-Chuan Jack Li. The journal was established in 1992 as Informatics in Primary Care, and was renamed to Journal of Innovation in Health Informatics in 2015, obtaining its current name in 2019.
It is an official publication of the British Computer Society and has a partnership agreement with the Faculty of Clinical Informatics.
Abstracting and indexing
The journal is abstracted and indexed in: Index Medicus/MEDLINE/PubMed, Scopus, and EBSCOhost.
References
External links
for BMJ Health & Care Informatics
for Journal of Innovation in Health Informatics (previous title)
Creative Commons-licensed journals
English-language journals
Biomedical informatics journals
Quarterly journals
Academic journals established in 1992
British Computer Society | BMJ Health & Care Informatics | [
"Biology"
] | 191 | [
"Bioinformatics",
"Biomedical informatics journals"
] |
49,352,276 | https://en.wikipedia.org/wiki/CXorf49 | CXorf49 is a protein, which in humans is encoded by the gene chromosome X open reading frame 49(CXorf49).
Gene
The CXorf49 gene has one alias CXorf49B. The recname A8MYA2 also refers to the protein coded by CXorf49 or CXorf49B.
CXorf49 is located on the X chromosome at Xq13.1. It is 3912 base pairs long and the gene sequence has 6 exons. CXorf49 has one protein coding transcript.
Protein
The protein has 514 amino acids and a molecular mass of 54.4 kDa. The isoelectric point is 9.3. Compared to other human proteins CXorf49 is glycine- and proline-rich, but the protein has lower levels of asparagine, isoleucine, tyrosine and threonine(Statistical Analysis of Protein Sequences, SAPS ).
Domains
The domain of unknown function, DUF4641, is almost the entire protein. It is 433 amino acids long, from amino acid 80 until amino acid number 512. DUF4641 is a part of pfam15483. The domain is proline- and arginine-rich, but DUF4641 has lower levels of isoleucine, tyrosine and threonine compared to other proteins in human (Analysis of Protein Sequences, SAPS ). DUF4641 has an unusual spacing between lysine residues and positive charged amino acids (Analysis of Protein Sequences, SAPS ).
Post-translation modifications
CXorf49 is predicted to have several post-translational sites. This include sites for N-acetyltransferase (NetAcet 1-), glycation of ε amino groups of lysines (NetGlycate 1.0), mucin type GalNAc O-glycosylation (NetOglyc 4.0), phosphorylation (NetPhos 2.0), sumoylation (SUMOplot Analysis Program) and O-ß-GlcNAc attachment(YinOYang WWW).
Subcellular localization
The CXorf49 protein has been predicted to be located in the cell nucleus (PSORT II ).
Expression
Promoter region
The promoter region of CXorf49 is located between base pair 71718051 and 71718785 on the minus strand of the X chromosome and it is 735 bp long (Genomatix’s ElDorado program). One of the most frequent transcription factor binding-sites in the promoter region are sites for Y-box binding factor.
Expression
Though expression of CXorf49 is very low in human cells, is it somewhat higher in connective tissues, testis and uterus(NCBI-Unigene ).
Interactions
The protein CXorf49 has not yet been shown to interact with other proteins (PSICQUIC).
CXorf49 is found to be one of the components of a small group of the HL-60 cell proteome that were most prone to form 4-Hydroxy-2-nonenal(HNE) adducts, upon exposure to nontoxic (10 μM) HNE concentrations, along with heat shock 60 kDa protein 1.
Homology
Using BLAST no orthologs for CXorf49 are found in single celled organisms, fungi or plants whose genomes have been sequenced. For multicellular organisms orthologs are found in mammals. The table below show a selection of the mammal orthologs. They are listed after time of divergence from human.
Phylogeny
CXorf49 has developed from aardvarks, to the human protein over 105.0 million years.
References
Genes
Proteins | CXorf49 | [
"Chemistry"
] | 796 | [
"Biomolecules by chemical classification",
"Proteins",
"Molecular biology"
] |
49,357,008 | https://en.wikipedia.org/wiki/Master%20regulator%20gene | In genetics, a master regulator gene is a regulator gene at the top of a gene regulation hierarchy, particularly in regulatory pathways related to cell fate and differentiation.
Examples
Most genes considered master regulators code for transcription factor proteins, which in turn alter the expression of downstream genes in the pathway. Canonical examples of master regulators include Oct-4 (also called POU5F1), SOX2, and NANOG, all transcription factors involved in maintaining pluripotency in stem cells. Master regulators involved in development and morphogenesis can also appear as oncogenes relevant to tumorigenesis and metastasis, as with the Twist transcription factor.
Other genes reported as master regulators code for SR proteins, which function as splicing factors, and some noncoding RNAs.
Criticism
The master regulator concept has been criticized for being a "simplified paradigm" that fails to account for the multifactorial influences on some cell fates.
References
Gene expression | Master regulator gene | [
"Chemistry",
"Biology"
] | 194 | [
"Gene expression",
"Molecular genetics",
"Cellular processes",
"Molecular biology",
"Biochemistry"
] |
49,359,017 | https://en.wikipedia.org/wiki/Zinc%20uptake%20regulator | The zinc uptake regulator (Zur) gene is a bacterial gene that codes for a transcription factor protein involved in zinc homeostasis. The protein is a member of the ferric uptake regulator family and binds zinc with high affinity. It typically functions as a repressor of zinc uptake proteins via binding to characteristic promoter DNA sequences in a dimer-of-dimers arrangement that creates strong cooperativity. Under conditions of zinc deficiency, the protein undergoes a conformational change that prevents DNA binding, thereby lifting the repression and causing zinc uptake genes such as ZinT and the ZnuABC zinc transporter to be expressed.
References
Prokaryote genes
Bacterial proteins
Transcription factors | Zinc uptake regulator | [
"Chemistry",
"Biology"
] | 143 | [
"Gene expression",
"Prokaryotes",
"Signal transduction",
"Induced stem cells",
"Prokaryote genes",
"Transcription factors"
] |
49,362,816 | https://en.wikipedia.org/wiki/Squire%27s%20theorem | In fluid dynamics, Squire's theorem states that of all the perturbations that may be applied to a shear flow (i.e. a velocity field of the form ), the perturbations which are least stable are two-dimensional, i.e. of the form , rather than the three-dimensional disturbances. This applies to incompressible flows which are governed by the Navier–Stokes equations. The theorem is named after Herbert Squire, who proved the theorem in 1933.
Squire's theorem allows many simplifications to be made in stability theory. If we want to decide whether a flow is unstable or not, it suffices to look at two-dimensional perturbations. These are governed by the Orr–Sommerfeld equation for viscous flow, and by Rayleigh's equation for inviscid flow.
References
Fluid dynamics | Squire's theorem | [
"Chemistry",
"Engineering"
] | 182 | [
"Piping",
"Chemical engineering",
"Fluid dynamics stubs",
"Fluid dynamics"
] |
45,093,519 | https://en.wikipedia.org/wiki/Refraction%20%28sound%29 | Refraction, in acoustics, comparable to the refraction of electromagnetic radiation, is the bending of sound propagation trajectories (rays) in inhomogeneous elastic media (gases, liquids, and solids) in which the wave velocity is a function of spatial coordinates. Bending of acoustic rays in layered inhomogeneous media occurs towards a layer with a smaller sound velocity. This effect is responsible for guided propagation of sound waves over long distances in the ocean and in the atmosphere.
In the atmosphere, vertical gradients of wind speed and temperature lead to refraction. The wind speed is usually increasing with height, which leads to a downward bending of the sound rays towards the ground. The same holds if the temperature is increasing with height (inversion). If the temperature is decreasing with height and the wind speed is low, sound rays are bent upwards.
See also
Atmospheric refraction
Deep sound channel
Sound speed gradient
Underwater acoustics
References
Further reading
P.M. Morse and K.U. Ingard, Theoretical Acoustics, Princeton University Press, 1986.
Acoustics | Refraction (sound) | [
"Physics"
] | 216 | [
"Classical mechanics",
"Acoustics"
] |
45,094,689 | https://en.wikipedia.org/wiki/Dipankar%20Das%20Sarma | Dipankar Das Sarma, popularly known as D.D. Sarma, is an Indian scientist and structural chemist, known for his researches in the fields of Solid State Chemistry, Spectroscopy, Condensed Matter Physics, Materials Science, and Nanoscience. He is a former MLS Chair Professor of Physics and Chairman of the Centre for Advanced Materials and the GAST Professor of Uppsala University, Sweden, A recipient of TWAS Physics Prize and the UNESCO Biennial Javed Husain Prize, Sarma was honored by the Council for Scientific and Industrial Research (CSIR), Government of India, in 1994, with the Shanti Swarup Bhatnagar Prize for Science and Technology.
Biography
Dipankar Das Sarma was born on 15 September 1955 in Kolkata, in West Bengal. He did a five-year integrated masters course in Physics from the Indian Institute of Technology, Kanpur in 1977 and enrolled for research at the Indian Institute of Science, (IISc) Bengaluru from where he secured his PhD in 1982 under the tutelage of renowned solid state chemist, C. N. R. Rao. He worked as a research associate at IISc for one year (1982–83), moved to Forschungszentrum Jülich, (Jülich Research Centre) Germany as a guest scientist in 1984 and returned to IISc as a lecturer in 1986. He stayed at IISc where he became the assistant professor in 1989, associate professor in 1993 and a professor in 1999. He remains a professor at Solid State and Structural Chemistry Unit at the institution. He also served as a visiting professor at the University of Tokyo (2001–02) and at the Istituto di Struttura della Materia, CNR at their Rome and Trieste centres in 2002.
Positions
Prof. Sarma holds a number of academic positions in India and abroad. Presently, he is J.N. Tata Chair Professor at the Solid State and Structural Chemistry Unit of the Indian Institute of Science. Besides heading the Solid State and Structural Chemistry Unit, he served as a guest professor at the Uppsala University, Sweden, a Distinguished Scientist at the CSIR-Network Institutes of Solar Energy, an honorary professor at the Jawaharlal Nehru Centre for Advanced Scientific Research, Bengaluru and a distinguished visiting professor at the Indian Association for the Cultivation of Science, Kolkata. He has also been an adjunct professor at the Tata Institute of Fundamental Research, Mumbai till October 2014 for a second term, the initial tenure being a six-year period from 2003 to 2009. He also held the position of the distinguished visiting professor at the Indian Association for the Cultivation of Science, Kolkata for two terms, from 2004 to 2006 and from 2009 to 2014. He has also been the adjunct professor of the Indian Institute of Science Education and Research, Kolkata (2007–09), a senior associate at the S.N. Bose National Centre for Basic Sciences (2003-2006), a member of faculty at the UGC-DAE Consortium for Scientific Research. and an MLS Chair professor of the Centre for Advanced Materials.
Prof. Sarma is a member of many professional bodies such as the Council of Raja Ramanna Centre for Advanced Tachnology, Indore, the Steering Committee for the Sophisticated Analytical Instrument Facilities (SAIF) programme of the Department of Science and Technology, Government of India and the Academic Advisory Committee of the Jawaharlal Nehru Centre for Advanced Scientific Research and is the coordinator of the group set up by the Government of India for India's collaboration with the International Centre for Synchrotron-Light for Experimental Science Applications in the Middle East (SESAME). He sits in various committees of the Department of Science and Technology and the Council for Scientific and Industrial Research. He holds the chair of the Research Council of the National Chemical Laboratory, Pune and has held the chairs of the Proposal Review Committee of Elettra Synchrotron Centre, Trieste and the Scientific Advisory Committee of the UGC-DAE Consortium for Scientific Research of which he has also served as a member of the governing body and the governing council.
D. D. Sarma is a former member of the Scientific Advisory Boards of CRANN and Trinity College Dublin, the Scientific Advisory Committees of the Inter University Consortium, Departamento de Asistencia Económica y Financiera, (DAEF) and the Inter University Accelerator Centre, New Delhi, the Academic Program Advisory Committee of the S.N. Bose National Centre for Basic Sciences, Kolkata and the councils of the Indian National Science Academy and the Raja Ramanna Centre for Advanced Technology, Indore. He has also served the scientific council of the Indo-French Centre for the Promotion of Advanced Research (IFCPAR), the general council of the Asia Pacific Centre for Theoretical Physics, Pohang, Korea, and the working group of the Department of Scientific and Industrial Research (DSIR) for the formulation of the 12th Five Year Plan, 2011 as a member.
Prof. Sarma is a Senior Editor of ACS Energy Letters and a member of the editorial boards of several peer reviewed journals such as J. Electron Spectroscopy and Related Phenomena, Solid State Communication, Indian Journal of Physics, and Surface and Interface Analysis. He has also been associated with Advances in Physical Chemistry, Pramana – Journal of Physics, Journal of Physical Chemistry, The Open Condensed Matter Physics Journal and Research Letters in Physical Chemistry, as an editorial board member. He is the Series Editor of Advances in Condensed Matter Science, the South Asian regional Editor for J. Experimental Nanoscience and the Associate Editor of Applied Physics. He is a reviewer for many scientific magazines and the American Physical Society selected him as an Outstanding Referee in 2009.
Legacy
Dipankar Das Sarma is credited with extensive research on nanomaterials and strongly correlated materials. His contribution is reported in discovering the existence of a new phase in solid state materials through high-energy spectroscopies and theory. His researches have been documented by way of several articles published in various peer reviewed journals. Google Scholar, an online repository of scientific articles has listed 656 of Sarma's articles and has accorded him an h-index of 43 (since 2015) and an i10-index of 167 (since 2015) and his articles have been cited over 20000 times. He holds many patents and Justia Patents has an online record of 16 of them.
Sarma is credited with the establishment of the Centre for Advanced Materials, a centre for advanced research on nanomaterials, smart materials, functional polymers, spintronics, strongly correlated electron systems, biomaterials and biology-inspired materials at the Indian Association for the Cultivation of Science. He has attended many national and international seminars and conferences where he has delivered plenary lectures and keynote addresses. He has also mentored many research students for their PhD theses.
Awards and recognitions
Sarma received the Sir J. C. Ghosh Medal in 1981 and the Young Scientist Medal from the Indian National Science Academy in 1983. UNESCO awarded him the Biennial Javed Hussain Prize in 1989 and a year later, in 1990, he received the medal of excellence from the Materials Research Society of India. The Government of India awarded him the Shanti Swarup Bhatnagar Prize in 1994 and he received the C. V. Raman Award in 2004. The year 2005 brought him three awards, the Hari Om Ashram Trust Award by the University Grants Commission, G. D. Birla Award and the Alumnus Award for Excellence in Research from the Indian Institute of Science. He received FICCI Award in 2006, TWAS Prize in Physics in 2007, the National Research Award in 2009 and H. K. Firodia Award in 2013. He was honored with the Knight of "The Order of the Star of Italy" in 2014, and an Honorary Doctorate from the Faculty of Science and Technology at Uppsala University in 2015.
Sarma, an elected Fellow of the Indian National Science Academy, delivered three INSA lectures in 2006, Dr. Jagdish Shankar Memorial Award Lecture, Professor R. P. Mitra Memorial Award Lecture and A. V. Rama Rao Foundation Lecture. Two years later, he delivered the INSA Kotcherlakota Rangadhama Rao Memorial Lecture in 2008. Some of the other lectures delivered by Das Sarma are CSIR Foundation Day Lecture, Distinguished Public Lecture, 8th Atma Ram Memorial Lecture and the Joy Kissen Memorial Lecture, CNR Rao Prize Lecture.
He is an elected Fellow of the American Physical Society, The Academy of Sciences for the Developing World, the National Academy of Sciences, India, and the Indian Academy of Sciences and holds the fellowship of the Asia-Pacific Academy of Materials (APAM). He is also a J. C. Bose National Fellow and Homi Bhabha Fellow.
See also
Indian Institute of Science
Indian Institute of Technology, Kanpur
Spintronics
Nanotechnology
References
1955 births
Living people
Indian chemical engineers
Recipients of the Shanti Swarup Bhatnagar Prize for Science and Technology
Scientists from Kolkata
Recipients of the Shanti Swarup Bhatnagar Award in Chemical Science
20th-century Indian inventors
Fellows of the Indian National Science Academy
Solid-state chemistry
Indian Institute of Science alumni
Academic staff of the Indian Institute of Science
20th-century Indian chemists
TWAS laureates
Fellows of the American Physical Society | Dipankar Das Sarma | [
"Physics",
"Chemistry",
"Materials_science"
] | 1,885 | [
"Condensed matter physics",
"nan",
"Solid-state chemistry"
] |
32,013,990 | https://en.wikipedia.org/wiki/Spin%20geometry | In mathematics, spin geometry is the area of differential geometry and topology where objects like spin manifolds and Dirac operators, and the various associated index theorems have come to play a fundamental role both in mathematics and in mathematical physics.
An important generalisation is the theory of symplectic Dirac operators in symplectic spin geometry and symplectic topology, which have become important fields of mathematical research.
See also
Contact geometry
Symplectic topology
Spinor
Spinor bundle
Spin manifold
Books
Differential topology
Differential geometry | Spin geometry | [
"Mathematics"
] | 103 | [
"Topology stubs",
"Topology",
"Differential topology"
] |
32,017,546 | https://en.wikipedia.org/wiki/%C3%89tienne%20Delaune | Étienne Delaune, Delaulne, or De Laune, (1518 or 1519) was a French goldsmith, medallist, draughtsman and engraver .
Life
He was born in Paris, or more probably at Orléans, in 1518. medallist, draughtsman He worked as a goldsmith in Paris in the 1550s. In 1552 he was appointed to the royal mint, where he would have produced metalwork designs. However, he left this position after six months following a dispute about wages. During his employment at the mint, Delaune had been able to build links with the French court and king. This allowed him to obtain royal commissions in particular for the design of intricately detailed royal armour, medals and other metalwork.
He commenced his career as an engraver of medals, and is said to have been helped by Benvenuto Cellini, who was at that time living in Paris. He afterwards engraved many prints after Raphael, and the Italian masters of Fontainebleau, and still more after the designs of his son Jean, with whom he passed the greater part of his life at Strassburg. His style was formed upon that of the Little Masters of Germany. He died at Strassburg around 1583.
Étienne Delaune was one of the most famous designers of goldsmithery of his time. There are six of his designs in the Louvre; two of them are for circular dishes representing the Histories of Moses and of Samson. His prints, which are generally small, are very numerous; they are executed entirely with the graver, with great dexterity of handling, and are very highly finished. He copied some of the prints of Marc Antonio with success. He usually marked his prints with the initial of his Christian name, S., or S. F., or S. fecit, but sometimes Stephanus, fecit.
Works
His works are described in Robert-Dumesnil's 'Peintre-Graveur,' vol. ix. The following are the principal:
A set of thirty Subjects from the Old Testament.
A set of eighteen Mythological Subjects; oval, very small.
The Twelve Months of the Year; circular.
Jupiter, Neptune, Mercury, and Ceres; four circular plates.
Four Subjects from Ancient History; oval.
The Four Monarchies; oval.
Four plates of Rural Occupations; oval.
The Three Graces.
David and Goliath; after Marc Antonio.
The Murder of the Innocents; after the same.
The Martyrdom of St. Felicitas; after the same.
The Rape of Helen; after the same.
The Brazen Serpent; after Jean Cousin. This is one of his largest prints.
References
Sources
External links
1510s births
1595 deaths
16th-century French engravers
French goldsmiths
Artists from Orléans
Metalsmiths from Paris
Material culture of royal courts | Étienne Delaune | [
"Engineering"
] | 576 | [
"Design engineering",
"Draughtsmen"
] |
32,018,036 | https://en.wikipedia.org/wiki/Mechanism%20of%20sonoluminescence | Sonoluminescence is a phenomenon that occurs when a small gas bubble is acoustically suspended and periodically driven in a liquid solution at ultrasonic frequencies, resulting in bubble collapse, cavitation, and light emission. The thermal energy that is released from the bubble collapse is so great that it can cause weak light emission. The mechanism of the light emission remains uncertain, but some of the current theories, which are categorized under either thermal or electrical processes, are Bremsstrahlung radiation, argon rectification hypothesis, and hot spot. Some researchers are beginning to favor thermal process explanations as temperature differences have consistently been observed with different methods of spectral analysis. In order to understand the light emission mechanism, it is important to know what is happening in the bubble's interior and at the bubble's surface.
Current competing theories
Prior to the early 1990s, the studies on different chemical and physical variables of sonoluminescence were all conducted using multi-bubble sonoluminescence (MBSL). This was a problem since all of the theories and bubble dynamics were based on single bubble sonoluminescence (SBSL) and researchers believed that the bubble oscillations of neighboring bubbles could affect each other. Single bubble sonoluminescence wasn't achieved until the early 1990s and allowed the study of the effects of various parameters on a single cavitating bubble. After many of the early theories were disproved, the remaining plausible theories can be classified into two different processes: electrical and thermal.
Single-bubble sonoluminescence (SBSL)
SBSL emits more light than MBSL due to fewer interactions between neighboring bubbles. Another advantage for SBSL is that a single bubble collapses without being affected by other surrounding bubbles, allowing more accurate studies on acoustic cavitation and sonoluminescence theories. Some exotic theories have been made, for example from Schwinger in 1992 who hinted the dynamical Casimir effect as a potential photon-emission process. Several theories say that the location of light emission is in the liquid instead of inside the bubble. Other SBSL theories explain that the emission of photons due to the high temperatures in the bubble are analogical to the hot spot theories of MBSL. Regarding the thermal emission a large variety of different processes are prevalent. Because temperatures are increasing from several hundred to many thousand kelvin during collapse, the processes can be molecular recombination, collision-induced emission, molecular emission, excimers, atomic recombination, radiative attachments of ions, neutral and ion Bremsstrahlung, or emission from confined electrons in voids. Which of these theories applies depends on accurate measurements and calculations of the temperature inside the bubble.
Multi-bubble sonoluminescence (MBSL)
Unlike single-bubble sonoluminescence, multi-bubble sonoluminescence is the creation of many oscillating and collapsing bubbles. Typically in MBSL, the light emission from each individual bubble is weaker than in SBSL because the neighboring bubbles can interact and affect each other. Because each neighboring bubble can interact with each other, it can make it more difficult to produce accurate studies and to characterize the properties of the collapsing bubble.
Bubble interior
One of the greatest obstacles in sonoluminescence research has been trying to obtain measurements of the interior of the bubble. Most measurements, like temperature and pressure, are indirectly measured using models and bubble dynamics.
Temperature
Some of the developed theories about the mechanism of SBSL result in prognoses for the peak temperature from 6000 K to 20,000 K. What they all have in common is, a) the interior of the bubble heats up and becomes at least as hot as that measured for MBSL, b) water vapor is the main temperature-limiting factor and c) the averaged temperature over the bubble does not rise higher than 10,000 K.
Bubble dynamics
These equations were made using five major assumptions, with four of them being common to all the equations:
The bubble remains spherical
The bubble contents obey the ideal gas law
The internal pressure remains uniform throughout the bubble
No evaporation or condensation occurs inside the bubble
The fifth assumption, which changes between each formulation, pertains to the thermodynamic behavior of the liquid surrounding the bubble. These assumptions severely limit the models when the pulsations are large and the wall velocities reach the speed of sound.
Keller–Miksis formulation
The Keller–Miksis formulation is an equation derived for the large, radial oscillations of a bubble trapped in a sound field. When the frequency of the sound field approaches the natural frequency of the bubble, it will result in large amplitude oscillations. The Keller–Miksis equation takes into account the viscosity, surface tension, incident sound wave, and acoustic radiation coming from the bubble, which was previously unaccounted for in Lauterborn's calculations. Lauterborn solved the equation that Plesset, et al. modified from Rayleigh's original analysis of large oscillating bubbles. Keller and Miksis obtained the following formula:
where is the radius of the bubble, the dots indicate first and second time derivatives, is the density of the liquid, is the speed of sound through the liquid, is the pressure on the liquid side of the bubble's interface, is time, and is the time-delayed driving pressure.
Prosperetti formulation
Prosperetti found a way to accurately determine the internal pressure of the bubble using the following equation.
where is the temperature, is the thermal conductivity of the gas, and is the radial distance.
Flynn's formulation
This formulation allows the study of the motions and the effects of heat conduction, shear viscosity, compressibility, and surface tension on small cavitation bubbles in liquids that are set into motion by an acoustic pressure field. The effect of vapor pressure on the cavitation bubble can also be determined using the interfacial temperature. The formulation is specifically designed to describe the motion of a bubble that expands to a maximum radius and then violently collapses or contracts. This set of equations was solved using an improved Euler method.
where is the radius of the bubble, the dots indicate first and second time derivatives, is the density of the liquid, is the speed of sound through the liquid, is the pressure on the liquid side of the bubble's interface, is time, and is the driving pressure.
Rayleigh–Plesset equation
The theory of bubble dynamics was started in 1917 by Lord Rayleigh during his work with the Royal Navy to investigate cavitation damage on ship propellers. Over several decades his work was refined and developed by Milton Plesset, Andrea Prosperetti, and others. The Rayleigh–Plesset equation is:
where is the bubble radius, is the second order derivative of the bubble radius with respect to time, is the first order derivative of the bubble radius with respect to time, is the density of the liquid, is the pressure in the gas (which is assumed to be uniform), is the background static pressure, is the sinusoidal driving pressure, is the viscosity of the liquid, and is the surface tension of the gas-liquid interface.
Bubble surface
The surface of a collapsing bubble like those seen in both SBSL and MBSL serves as a boundary layer between the liquid and vapor phases of the solution.
Generation
MBSL has been observed in many different solutions under a variety of conditions. Unfortunately it is more difficult to study as the bubble cloud is uneven and can contain a wide range of pressures and temperatures. SBSL is easier to study due to the predictable nature of the bubble. This bubble is sustained in a standing acoustic wave of moderate pressure, approximately 1.5 atm. Since cavitation does not normally occur at these pressures the bubble may be seeded through several techniques:
Transient boiling through short current pulse in nichrome wire.
A small jet of water perturbs the surface to introduce air bubbles.
A rapidly formed vapor cavity via focused laser pulse.
The standing acoustic wave, which contains pressure antinodes at the center of the containment vessel, causes the bubbles to quickly coalesce into a single radially oscillating bubble.
Collapse
Once a single bubble is stabilized in the pressure antinode of the standing wave, it can be made to emit pulses of light by driving the bubble into highly nonlinear oscillations. This is done by the increasing pressure of the acoustic wave to disrupt the steady, linear growth of the bubble which cause the bubble to collapse in a runaway reaction that only reverts due to the high pressures inside the bubble at its minimum radius.
Afterbounces
The collapsed bubble expands due to high internal pressure and experiences a diminishing effect until the high pressure antinode returns to the center of the vessel. The bubble continues to occupy more or less the same space due to the acoustic radiation force, the Bjerknes force, and the buoyancy force of the bubble.
Surface chemistry
The effect that different chemicals present in solution have to the velocity of the collapsing bubble has recently been studied. Nonvolatile liquids such as sulfuric and phosphoric acid have been shown to produce flashes of light several nanoseconds in duration with a much slower bubble wall velocity, and producing several thousand-fold greater light emission. This effect is probably masked in SBSL in aqueous solutions by the absorption of light by water molecules and contaminants.
Surface tension
It can be inferred from these results that the difference in surface tension between these different compounds is the source of different spectra emitted and the time scales in which emission occur.
Light emission
The inertia of a collapsing bubble generates high pressures and temperatures capable of ionizing a small fraction of the noble gas within the volume of the bubble. This small fraction of ionized gas is transparent and allows for volume emission to be detected. Free electrons from the ionized noble gas begin to interact with other neutral atoms causing thermal bremsstrahlung radiation. Surface emission emits a more intense flash of light with a longer duration and is dependent on wavelength. Experimental data suggest that only volume emission occurs in the case of sonoluminescence. As the sound wave reaches a low energy trough, the bubble expands and electrons are able to recombine with free ions and halt light emission. Light pulse time is dependent on the ionization energy of the noble gas with argon having a light pulse of 160 picoseconds.
Electrical processes
In 1937, the explanations for the light emission have favored electrical discharges. The first ideas have been about the charge separation in cavitation bubbles, which have been seen as spherical capacitors with charges at the center and the wall.
At the collapse, the capacitance decreases and voltage increases until electric breakdown occurs. A further suggestion was a charge separation by enhancing charge fluctuations on the bubble wall, however, a breakdown should take place during the expansion phase of the bubble dynamics.
These discharge theories have to assume that the emitting bubble undergoes an asymmetric collapse, because a symmetric charge distribution cannot radiate light.
Thermal processes
Because the bubble collapse occurs within microseconds, the hot spot theory states that the thermal energy results from an adiabatic bubble collapse. In 1950 it was assumed that the bubble internal temperatures were as high as 10,000 K at the collapse of a spherical symmetric bubble. In the 1990s, sonoluminescence spectra were used by Suslick to measure effective emission temperatures in bubble clouds (multibubble sonoluminescence) of 5000 K, and more recently temperatures as high as 20,000 K in single bubble cavitation.
Bubble shape stability
The limit for the ambient size of the bubble is set by the appearance of instabilities in the shape of the oscillating bubble.
The shape stability thresholds depend on changes in the radial dynamics, caused by different liquid viscosities or driving frequencies. If the frequency is decreased, the parametric instability is suppressed as the stabilizing influence of viscosity can appear longer to suppress perturbations. However, the collapses of low-frequency-driven bubbles favor an earlier onset of the Rayleigh-Taylor instability. Larger bubbles can be stabilized to show sonoluminescence when not too high forcing pressures are applied. At low-frequency the water vapor becomes more important. The bubbles can be stabilized by cooling the fluid, whereas more light is emitted.
See also
Bubble fusion
References
Luminescence | Mechanism of sonoluminescence | [
"Chemistry"
] | 2,558 | [
"Luminescence",
"Molecular physics"
] |
32,018,124 | https://en.wikipedia.org/wiki/Amino%20acid%20kinase | In molecular biology, the amino acid kinase domain is a protein domain. It is found in protein kinases with various specificities, including the aspartate, glutamate and uridylate kinase families. In prokaryotes and plants the synthesis of the essential amino acids lysine and threonine is predominantly regulated by feed-back inhibition of aspartate kinase (AK) and dihydrodipicolinate synthase (DHPS). In Escherichia coli, thrA, metLM, and lysC encode aspartokinase isozymes that show feedback inhibition by threonine, methionine, and lysine, respectively. The lysine-sensitive isoenzyme of aspartate kinase from spinach leaves has a subunit composition of 4 large and 4 small subunits.
In plants although the control of carbon fixation and nitrogen assimilation has been studied in detail, relatively little is known about the regulation of carbon and nitrogen flow into amino acids. The metabolic regulation of expression of an Arabidopsis thaliana aspartate kinase/homoserine dehydrogenase (AK/HSD) gene, which encodes two linked key enzymes in the biosynthetic pathway of aspartate family amino acids has been studied. The conversion of aspartate into either the storage amino acid asparagine or aspartate family amino acids may be subject to a coordinated, reciprocal metabolic control, and this biochemical branch point is a part of a larger, coordinated regulatory mechanism of nitrogen and carbon storage and utilization.
References
Protein families | Amino acid kinase | [
"Biology"
] | 333 | [
"Protein families",
"Protein classification"
] |
32,020,153 | https://en.wikipedia.org/wiki/Massive%20parallel%20sequencing | Massive parallel sequencing or massively parallel sequencing is any of several high-throughput approaches to DNA sequencing using the concept of massively parallel processing; it is also called next-generation sequencing (NGS) or second-generation sequencing. Some of these technologies emerged between 1993 and 1998 and have been commercially available since 2005. These technologies use miniaturized and parallelized platforms for sequencing of 1 million to 43 billion short reads (50 to 400 bases each) per instrument run.
Many NGS platforms differ in engineering configurations and sequencing chemistry. They share the technical paradigm of massive parallel sequencing via spatially separated, clonally amplified DNA templates or single DNA molecules in a flow cell. This design is very different from that of Sanger sequencing—also known as capillary sequencing or first-generation sequencing—which is based on electrophoretic separation of chain-termination products produced in individual sequencing reactions. This methodology allows sequencing to be completed on a larger scale.
NGS platforms
DNA sequencing with commercially available NGS platforms is generally conducted with the following steps. First, DNA sequencing libraries are generated by clonal amplification by PCR in vitro. Second, the DNA is sequenced by synthesis, such that the DNA sequence is determined by the addition of nucleotides to the complementary strand rather than through chain-termination chemistry. Third, the spatially segregated, amplified DNA templates are sequenced simultaneously in a massively parallel fashion without the requirement for a physical separation step. These steps are followed in most NGS platforms, but each utilizes a different strategy.
NGS parallelization of the sequencing reactions generates hundreds of megabases to gigabases of nucleotide sequence reads in a single instrument run. This has enabled a drastic increase in available sequence data and fundamentally changed genome sequencing approaches in the biomedical sciences.
Newly emerging NGS technologies and instruments have further contributed to a significant decrease in the cost of sequencing nearing the mark of $1000 per genome sequencing.
As of 2014, massively parallel sequencing platforms are commercially available and their features are summarized in the table. As the pace of NGS technologies is advancing rapidly, technical specifications and pricing are in flux.
Run times and gigabase (Gb) output per run for single-end sequencing are noted. Run times and outputs approximately double when performing paired-end sequencing.
‡Average read lengths for the Roche 454 and Helicos Biosciences platforms.
Template preparation methods for NGS
Two methods are used in preparing templates for NGS reactions: amplified templates originating from single DNA molecules, and single DNA molecule templates.
For imaging systems which cannot detect single fluorescence events, amplification of DNA templates is required. The three most common amplification methods are emulsion PCR (emPCR), rolling circle and solid-phase amplification. The final distribution of templates can be spatially random or on a grid.
Emulsion PCR
In emulsion PCR methods, a DNA library is first generated through random fragmentation of genomic DNA. Single-stranded DNA fragments (templates) are attached to the surface of beads with adaptors or linkers, and one bead is attached to a single DNA fragment from the DNA library. The surface of the beads contains oligonucleotide probes with sequences that are complementary to the adaptors binding the DNA fragments. The beads are then compartmentalized into water-oil emulsion droplets. In the aqueous water-oil emulsion, each of the droplets capturing one bead is a PCR microreactor that produces amplified copies of the single DNA template.
Gridded rolling circle nanoballs
Amplification of a population of single DNA molecules by rolling circle amplification in solution is followed by capture on a grid of spots sized to be smaller than the DNAs to be immobilized. Second-generation sequencing technologies like MGI Tech's DNBSEQ or Element Biosciences' AVITI use this approach for the preparation of the sample on the flow cell that is then imaged cycle by cycle.
DNA colony generation (Bridge amplification)
Forward and reverse primers are covalently attached at high-density to the slide in a flow cell. The ratio of the primers to the template on the support defines the surface density of the amplified clusters. The flow cell is exposed to reagents for polymerase-based extension, and priming occurs as the free/distal end of a ligated fragment "bridges" to a complementary oligo on the surface. Repeated denaturation and extension results in localized amplification of DNA fragments in millions of separate locations across the flow cell surface. Solid-phase amplification produces 100–200 million spatially separated template clusters, providing free ends to which a universal sequencing primer is then hybridized to initiate the sequencing reaction. This technology was filed for a patent in 1997 from Glaxo-Welcome's Geneva Biomedical Research Institute (GBRI), by Pascal Mayer, Eric Kawashima, and Laurent Farinelli, and was publicly presented for the first time in 1998. In 1994 Chris Adams and Steve Kron filed a patent on a similar, but non-clonal, surface amplification method, named “bridge amplification” adapted for clonal amplification in 1997 by Church and Mitra.
Single-molecule templates
Protocols requiring DNA amplification are often cumbersome to implement and may introduce sequencing errors. The preparation of single-molecule templates is more straightforward and does not require PCR, which can introduce errors in the amplified templates. AT-rich and GC-rich target sequences often show amplification bias, which results in their underrepresentation in genome alignments and assemblies.
Single molecule templates are usually immobilized on solid supports using one of at least three different approaches. In the first approach, spatially distributed individual primer molecules are covalently attached to the solid support. The template, which is prepared by randomly fragmenting the starting material into small sizes (for example,~200–250 bp) and adding common adapters to the fragment ends, is then hybridized to the immobilized primer. In the second approach, spatially distributed single-molecule templates are covalently attached to the solid support by priming and extending single-stranded, single-molecule templates from immobilized primers. A common primer is then hybridized to the template.
In either approach, DNA polymerase can bind to the immobilized primed template configuration to initiate the NGS reaction. Both of the above approaches are used by Helicos BioSciences. In a third approach, spatially distributed single polymerase molecules
are attached to the solid support, to which a primed template molecule is bound. This approach is used by Pacific Biosciences. Larger DNA molecules (up to tens of thousands of base pairs) can be used with this technique and, unlike the first two approaches, the third approach can be used with real-time methods, resulting in potentially longer read lengths.
Sequencing approaches
Sequencing by synthesis
The objective for sequential sequencing by synthesis (SBS) is to determine the sequencing of a DNA sample by detecting the incorporation of a nucleotide by a DNA polymerase. An engineered polymerase is used to synthesize a copy of a single strand of DNA and the incorporation of each nucleotide is monitored. The principle of sequencing by synthesis was first described in 1993 with improvements published some years later. The key parts are highly similar for all embodiments of SBS and include (1) amplification of DNA to enhance the subsequent signal and to attach the DNA to be sequenced to a solid support, (2) generation of single stranded DNA on the solid support, (3) incorporation of nucleotides using an engineered polymerase and (4) detection of the incorporation of nucleotide. Then steps 3-4 are repeated and the sequence is assembled from the signals obtained in step 4. This principle of sequencing-by-synthesis has been used for almost all massive parallel sequencing instruments, including 454, PacBio, IonTorrent, Illumina and MGI.
Pyrosequencing
The principle of Pyrosequencing was first described in 1993 by combining a solid support with an engineered DNA polymerase lacking 3´to 5´exonuclease activity (proof-reading) and luminescence real-time detection using the firefly luciferase. All the key concepts of sequencing by synthesis were introduced, including (1) amplification of DNA to enhance the subsequent signal and attach the DNA to be sequenced (template) to a solid support, (2) generation of single stranded DNA on the solid support (3) incorporation of nucleotides using an engineered polymerase and (4) detection of the incorporated nucleotide by light detection in real-time. In a follow-up article, the concept was further developed and in 1998, an article was published in which the authors showed that non-incorporated nucleotides could be removed with a fourth enzyme (apyrase) allowing sequencing by synthesis to be performed without the need for washing away non-incorporated nucleotides.
Sequencing by reversible terminator chemistry
This approach uses reversible terminator-bound dNTPs in a cyclic method that comprises nucleotide incorporation, fluorescence imaging and cleavage.
A fluorescently-labeled terminator is imaged as each dNTP is added and then cleaved to allow incorporation of the next base.
These nucleotides are chemically blocked such that each incorporation is a unique event. An imaging step follows each base incorporation step, then the blocked group is chemically removed to prepare each strand for the next incorporation by DNA polymerase. This series of steps continues for a specific number of cycles, as determined by user-defined instrument settings. The 3' blocking groups were originally conceived as either enzymatic or chemical reversal The chemical method has been the basis for the Solexa and Illumina machines.
Sequencing by reversible terminator chemistry can be a four-colour cycle such as used by Illumina/Solexa, or a one-colour cycle such as used by Helicos BioSciences.
Helicos BioSciences used “virtual Terminators”, which are unblocked terminators with a second nucleoside analogue that acts as an inhibitor. These terminators have the appropriate modifications for terminating or inhibiting groups so that DNA synthesis is terminated after a single base addition.
Sequencing-by-ligation mediated by ligase enzymes
In this approach, the sequence extension reaction is not carried out by polymerases but rather by DNA ligase and either one-base-encoded probes or two-base-encoded probes. In its simplest form, a fluorescently labelled probe hybridizes to its complementary sequence adjacent to the primed template. DNA ligase is then added to join the dye-labelled probe to the primer. Non-ligated probes are washed away, followed by fluorescence imaging to determine the identity of the ligated probe.
The cycle can be repeated either by using cleavable probes to remove the fluorescent dye and regenerate a 5′-PO4 group for subsequent ligation cycles (chained ligation) or by removing and hybridizing a new primer to the template (unchained ligation).
Phospholinked Fluorescent Nucleotides or Real-time sequencing
Pacific Biosciences is currently leading this method.
The method of real-time sequencing involves imaging the continuous incorporation of dye-labelled nucleotides during DNA synthesis: single DNA polymerase molecules
are attached to the bottom surface of individual zero-mode waveguide detectors (Zmw detectors) that can obtain sequence information while phospholinked nucleotides are being incorporated into the growing primer strand.
Pacific Biosciences uses a unique DNA polymerase which better incorporates phospholinked nucleotides and enables the resequencing of closed circular templates.
While single-read accuracy is 87%, consensus accuracy has been demonstrated at 99.999% with multi-kilobase read lengths. In 2015, Pacific Biosciences released a new sequencing instrument called the Sequel System, which increases capacity approximately 6.5-fold.
See also
Clinical metagenomic sequencing
First-generation sequencing
Third-generation sequencing
RNA Velocity
References
DNA sequencing methods | Massive parallel sequencing | [
"Biology"
] | 2,558 | [
"Genetics techniques",
"DNA sequencing methods",
"DNA sequencing"
] |
33,596,365 | https://en.wikipedia.org/wiki/Astron%20%28fusion%20reactor%29 | The Astron is a type of fusion power device pioneered by Nicholas Christofilos and built at the Lawrence Livermore National Laboratory during the 1960s and 70s. Astron used a unique confinement system that avoided several of the problems found in contemporary designs like the stellarator and magnetic mirror. Development was greatly slowed by a series of changes to the design that were made with limited oversight, leading to a review committee being set up to oversee further development. The Astron was unable to meet the performance goals set for it by the committee; funding was cancelled in 1972 and development wound down in 1973. Work on similar designs appears to have demonstrated a theoretical problem in the very design that suggests it could never be used for practical generation.
History
Strong focusing
Christofilos is best known for independently inventing the concept of strong focusing, a feature used in particle accelerators. He had first started work along these lines in the late 1940s while running an elevator installation company in Greece, and in 1948 he wrote a letter to what was then the University of California's Radiation Laboratory at Berkeley outlining several ideas on accelerator focusing. When they returned his letter pointing out several problems, he solved these and wrote them again. This second letter was ignored. In 1950 Christofilos filed a patent application, which was granted in 1956 as US Patent 2,736,799.
Around the same time, Ernest Courant, Milton Stanley Livingston, and Hartland Snyder of Brookhaven National Laboratory were considering the same problem and devised the same solution, writing about it in the 1 December 1952 issue of Physical Review. When he saw the paper, Christofilos arranged a trip to the US, arriving two months later. Making his way to Brookhaven, he angrily accused them of stealing the idea from his patent. He also met with members of the Atomic Energy Commission, and after meeting with his attorneys, they paid him $10,000 for the patent.
Astron proposal
With the patent purchase came some fame and enough money that Christofilos was able to enter the US physics world. In April 1953 he attended a meeting of Project Sherwood, and presented another idea he had been working on in Greece, the Astron.
The basic idea was to inject high-energy electrons into a magnetic mirror (the "tank"). The electrons would be captured in the mirror, and build up a layer of current near the outside surface of the tank volume, which he called the "E-layer". The E-layer would itself produce a powerful magnetic field as it built up, and once the current reached a critical density, the fields would "reverse", and fold into a new configuration of closed lines that formed a continuous confinement area. Once the E-layer had successfully formed, fusion fuel would be injected into the area inside it, and heated by interactions with the E-layer to bring it to fusion temperatures.
This arrangement solved one of the main problems with the basic magnetic mirror concept, which had open field lines at the ends. Fuel could follow these lines right out of the reactor. Mirrors thus naturally leaked plasma, although designers believed they could address this problem by operating the machines at very high temperatures. In practice the leakage proved to be even higher than basic theory suggested, and never operated at the levels they hoped to achieve.
At the time, Sherwood was still secret, which presented problems when he first outlined the concept. Prior to his taking the stage, the formulas from the previous session on the blackboard had been carefully erased. As he filled the blackboard with his own equations, someone helpfully showed him the buttons that would raise it and reveal a fresh one underneath. This one had not been erased and led to a rushed effort to prevent any sensitive material leaking. Looking to avoid a repeat event, Christofilos was given a job at Brookhaven, where he could continue working on the Astron theory.
Astron testing
In 1956, Christofilos finally received his security clearance, and he immediately moved to what was now the Lawrence Livermore National Laboratory (LLNL) to start work on the Astron concept. After two years enough progress had been made that he was able to present the idea at the 1958 Atoms for Peace conference in Geneva, along with a model of the system they proposed to build. This consisted of two main parts, the magnetic bottle where the plasma would be held, and a particle accelerator that provided the relativistic electrons.
In spite of his success, Christofilos was always an outsider at the lab. Time reported that "He still has no degree in physics, and his Greek accent, Greek volubility and love of passionate argument keep him an outsider." This led to friction within the physics establishment, and early calls for termination of the Astron program. A 1963 review of the entire Project Sherwood effort led to formal calls for cancellation. However, the program had backers within the management of the controlled fusion program, notably Glenn Seaborg and John S. Foster, both with strong ties to LLNL. Foster, in particular, was concerned about groups in Washington dictating development to the labs. After considerable argument, it was decided that the program would be allowed to continue, but would need to demonstrate field reversal by 1965.
By 1963 the team had designed and built a new type of linear induction accelerator with the required properties. The accelerator design led to interest as a particle beam weapon studied under Project Seesaw. However, during construction the team realized that the electrons were free to travel back into the accelerator area. Christofilos solved this by introducing resistor wires that slightly slowed the electrons after entering the tank, so they no longer possessed the energy needed to flow back out.
After some work ironing out bugs, the first results were published in June 1964. The accelerator worked, operating at 4 MeV and 120 amps, and a stable E-layer was confirmed, albeit generating only 2 A/cm of current, just 0.05% of the diamagnetic field required to reverse the field. Work continued to meet the 1965 goal of reversal, but ultimately failed. However, the electron layer was stable, so the Herb-Allison committee recommended it continue to the next milepost.
By 1967 this had been improved to 6%, but was still a long way from the stable E-layer the device needed to achieve. In 1968, Christofilos and T. Kenneth Fowler wrote a report asking for a more powerful accelerator, and upgrades to the tank.
Scrutiny
Funds for the upgrades were eventually granted, but only at the cost of direct oversight by an Ad Hoc Panel created by the AEC. By this point the "conventional" designs, the stellarator and magnetic mirror, had long been working on real-world plasmas and were slowly increasing the pressures and temperatures. Astron, on the other hand, was still a long way from building its first useful E-layer, a prerequisite for plasma experiments.
The Ad Hoc Panel returned a negative report, complaining that far too much effort had been put into operational issues like accelerator performance, with little or no effort into theoretical studies on whether or not the plasma would ever be stable even if an E-layer could be formed. Moreover, the panel pointed out that no one had seriously studied whether or not an operating and stable Astron would require more power to operate than it would release. This was a serious concern in Astron, because its relativistic electrons would radiate away large amounts of power due to electron synchrotron radiation.
Christofilos had already considered this and suggested that an operational design would use protons in place of the electrons, and would not suffer from the same level of energy losses. However, no such accelerator existed at the time, and the panel was highly skeptical that it would be simple to build.
Upgrade
The upgrades to Astron went ahead and started operation in 1969. During this period, following the advice of the Panel, the theoretical divisions at LLNL started taking a much more serious look at the concept. Building computer models of the system, they first attacked the problem of "stacking", that individual pulses of electrons from the accelerator did not build up in the E-layer as expected. Bruce Langdon demonstrated that stacking simply would not work.
However, a suggestion by Fowler proved to save the Astron from this problem. He had noted that adding a second magnetic field running down the centre of the tank would reduce the amount of external field needed to create the E-layer. Christofilos went ahead with this change and started testing in 1971; this demonstrated greatly improved performance both with the reduction in current and success in trapping the electrons. This also allowed two pulses to be stacked, raising the field to 15% diamagnetic strength.
While Astron worked towards multiple pulses, a team at Cornell University had been working on a similar design. However, this Relativistic Electron Coil Experiment (RECE) used a single long pulse of electrons rather than the stacking concept. Late in 1971 they announced they had achieved complete field reversal. Christofilos was unimpressed; this design would not be useful for a steady state fusion generator, only by continually adding pulses could the machine maintain itself.
Cancellation
Faced with the continued problems with Astron, and the seeming ease that the RECE team had managed to reach the goals they had originally suggested in 1968, a second Ad Hoc Panel published a scathing report. Among the problems they noted that the Astron team had been looking "for ingenious ways to avoid or circumvent difficulties rather than to understand them." Roy Gould, head of the AEC's controlled fusion program, was specific in allowing the Astron project to continue, but only if it met a series of goals on a specific timeline.
When Robert Hirsch took over the AEC's controlled fusion arm in 1972, he instituted a sweeping review to classify the approaches under study and eliminate duplication and low-payoff projects. Given the exciting results with the tokamak released in 1968, Hirsch favoured a program with relatively few projects each given much larger budgets. Many programs like Astron simply didn't appear to have any near-term payoff, and Hirsch was keen to cancel them.
On 24 September 1972, Christofilos met with James Schlesinger of the AEC, but no record of the meeting remains. After a long day, he went to a local Holiday Inn to save a long commute home. That night he suffered a massive heart attack and died.
Richard Briggs took over direction of the project until its planned shutdown date in June 1973. Under his direction, Astron returned to study of the new stabilizing field introduced by Fowler, and using single larger pulses the device hit 50% diamagnetic strength, much greater than Christofilos' efforts with pulse chains. Their final report stated that "buildup of the E-layer by multiple-pulse injection was generally unsuccessful" and noted that at the time of the shutdown they still did not understand what physics problem was limiting the buildup.
After Astron
Although Astron shut down, work continued with RECE at Cornell for some time. As part of their work, the team attempted to make the switch from electrons to protons. However, as some suspected, the "P-layer" proved difficult to build, and field reversal with protons was never achieved. The last version of this effort, FIREX, shut down in 2003, having demonstrated what appears to be a purely theoretical reason why the Astron concept will never work.
The relativistic electron ring also played a part in the bumpy torus design. This was another attempt to "plug the ends" of mirrors, by linking a number of mirrors end-to-end to form a torus. Electrons were driven to high energies not through direct injection, but external microwave-driven electron-cyclotron heating (ECH).
Description
The Astron device consisted of two sections, the linear accelerator and the magnetic mirror "tank". These were constructed at right angles, with the accelerator's output firing into the side of the tank at one end.
The tank was a relatively simple example of the magnetic mirror concept, consisting largely of a long solenoid with additional windings at both ends to increase the magnetic field in those regions and form the mirror. In a simple mirror the ions in the fuel plasma were injected at an angle so they could not simply flow right out of the ends where the field was roughly linear. However, there was an annular region on either end where ions of the right energy could escape, and various calculations demonstrated the rate would be fairly high.
By injecting electrons into the mirror before the fuel, the E-layer would create a second magnetic field that would cause the annular areas to fold back into the center of the tank. The resulting field was shaped like a tube, and very similar to the Field-Reversed Configuration, or FRC. The main difference between these devices is the way the field reversal is achieved; with the E-layer in the Astron, and by currents in the plasma for the FRC. Like the classic mirror, Astron injected the electrons into the mirror at a slight angle to ensure they would circulate into the center of the mirror.
Today, the Astron is often considered a sub-class of the FRC concept.
References
Citations
Bibliography
Fusion reactors
Lawrence Livermore National Laboratory | Astron (fusion reactor) | [
"Chemistry"
] | 2,741 | [
"Nuclear fusion",
"Fusion reactors"
] |
33,599,550 | https://en.wikipedia.org/wiki/Nanoreactor | Nanoreactors are a form of chemical reactor that are particularly in the disciplines of nanotechnology and nanobiotechnology. These special reactors are crucial in maintaining a working nanofoundry; which is essentially a foundry that manufactures products on a nanotechnological scale.
Summary
General information
The term nanoreactor refers to an isolated system on the nanometer scale that is used to run chemical reactions in an environment that differs drastically from a reaction in bulk solution. The synthesis and analysis of these nanoreactors is a highly interdisciplinary subject, spanning from chemistry and physics to biology and materials science. These systems can be synthetic, such as nanopores and hollow nanoparticles, or they can be biological systems, including protein pores and channels. Generally, the effect of confinement provided by these nanoreactors results in novel chemistry. This field has only begun to receive significant attention in the last two decades, and more work is constantly being published as nanoreactors become more sophisticated and begin to show promise for industrial applications.
Researchers in the Netherlands have succeeded in building nanoreactors that can perform one-pot multistep reactions - the next step towards artificial cell-like devices in addition for applications involving the screening and diagnosis of a disease or illness. A biochemical nanoreactor is created simply by unwrapping a biological virus through scientific methods, eliminating its harmful contents, and re-assembling its protein coat around a single molecule of enzyme. The kinetic isotope effect is trapped in a single molecule within a membrane-based nanoreactor. This is a phenomenon that has been found by researchers in the United Kingdom during experiments done in September 2010. The kinetic isotope effect, where the rate of a reaction is influenced by the presence of an isotopic atom in solution, is an important principle for elucidating reaction mechanisms. This recent finding could open up new methods to study chemical reactions. They may even aid in the process of creating new (and even more powerful) nanoreactors.
Using nanocrystals, a scalable and inexpensive process can ultimately create nanoreactors. Researchers at the Lawrence Berkeley National Laboratory in Berkeley have the ability to take advantage of the large difference in select components to create these nanocrystals and nanoreactors. Nanocrystals are easier to use and less expensive than methods that employ sacrificial templates in the creation process of hollow particles. Catalyst particles are separated into shells in order to prevent particle aggregation. Selective entry into the catalysis chamber reduces the likelihood of desired products undergoing secondary reactions.
Nanoreactors can also be built by controlling the positioning of two different enzymes in the central water reservoir or the plastic membrane of synthetic nanoscopic bubbles. Once the third enzyme is added into the surrounding solution, it becomes possible for three different enzymatic reactions to occur at once without interfering with each other (resulting in a "one-pot" reaction). The potential for nanoreactors can be demonstrated by binding the enzyme horseradish peroxidase into the membrane itself; trapping the enzyme glucose oxidase. The surrounding solution would end up containing the enzyme lipase B with the glucose molecules containing four acetyl groups as the substrate. The resulting glucose would cross the membrane, become oxidized, and the horseradish peroxidase would convert the sample substrate ABTS (2,2’-azinobis(3-ethylbenzthiazoline-6-sulfonic acid)) into its radical cation.
Abilities
Nanoreactors can also be used to emulsify water, create hydrofuels (which essentially blends 15% water into the refined diesel product), play a helpful role in the chemical industry by allowing multiple streams of raw materials to exists in a single nanoreactor, manufacture personal care products (i.e., lotions, pharmaceutical creams, shampoos, conditioners, shower gels, deodorants), and improve the food and beverage industries (by processing sauces, purées, cooking bases for soup, emulsifying non-alcoholic beverages, and salad dressings).
Personal care goods can be enhanced by companies feeding multiple phases of material, using a mixing device with water, and creating instant emulsions. These emulsions would come with smaller particles, are expected to have a longer shelf life and an give off an enhanced appearance when sold at retailers. The needs of the food and beverage industry can result in lower processing costs, more space, better efficiency, and lower equipment costs. This may bring down the cost of food and beverages for consumers; even alcoholic beverages that are subject to hidden sin taxes.
Hydrofuel can be used to move heavy duty transports, trains, earth-moving equipment (including bulldozers), in addition to providing fuel to most boats and ships. Reduced pollution and increased fuel efficiency may come out of nanoreactor-produced hydrofuel. The increased usage of renewable energy may also help to improve the world's environment thanks to nanoreactors.
Applications
Roy, Skinner, et al. studied the dynamics of water in self-assembled gemini surfactants in 2014. This work illustrates not only the utility of nano-scale materials for chemical reactions, but also the complexity that is required to study the effects. The team utilized spectroscopic techniques and molecular dynamics simulations to determine that within the nanoporous structures, the dynamics of water in the gyroid phase is an order of magnitude slower than in the bulk water. This result arises from the difference in curvature at the interfaces of the normal gyroid. When compared with water confined in a reverse spherical micelle of a sulfonate surfactant, the water exhibited faster dynamics. This complex behavior was postulated to have implications for future work in ion transport.
Carbon nanotubes have been a popular area of research, and specifically, single-walled carbon nanotubes provide unique surfaces for chemistry. Li, G and Fu, C et al. report on large changes to the Raman spectra by encapsulating sulfur in these single-walled carbon nanotubes. In an example of how confinement to such small spaces influences chemistry, the authors theorize that the changes to the Raman spectra can be attributed to van der Waals interactions of the sulfur with the walls of the nanotubes. These effects are highly sensitive to the size of the confinement chamber, as the van der Waals interactions were not significant for larger diameter single-walled nanotubes. The authors suggest that confinement within the single-walled nanotubes allows S2 molecules to undergo polymerization to linear diradicals.
Nanoreactors are also being applied to biological spaces. In a study by Tagliazucchi and Szleifer, they study the binding of proteins to ligands inside of both long nanochannels and short nanopores. Inside these confined spaces, the ligands are attached to the walls by polymeric tethers. This technology has already seen applications as sensors that measure concentrations of proteins in solution. This study developed a theory to model how the proteins bind under these highly confined conditions to inform the design of these sensors.
References
Nanotechnology | Nanoreactor | [
"Materials_science",
"Engineering"
] | 1,460 | [
"Nanotechnology",
"Materials science"
] |
33,603,101 | https://en.wikipedia.org/wiki/Centre%20for%20Ships%20and%20Ocean%20Structures | The Centre for Ships and Ocean Structures (CeSOS) is a research centre located at the Marine Technology Centre in Trondheim, Norway. The research centre's goal is to create fundamental knowledge about the design and operation of ships and ocean structures. The centre has been active since 2002, when it was established as a Centre of Excellence (CoE) by the Research Council of Norway and the Norwegian University of Science and Technology (NTNU). Although the financing period by the Research Council of Norway finished in 2012, research activities are still ongoing in 2013 and 2014, financed by external means.
Research areas
Research at CeSOS focuses on the creation of fundamental knowledge about the design and operation of future ships and ocean structures, using analytical, numerical and experimental studies. The centre has furthermore made it its responsibility to integrate marine hydrodynamics, structural mechanics and automatic control in its research. The knowledge that is being produced by this kind of research is vital, both now and in the future, for the design of safe, cost effective and environmentally friendly structures as well as in the planning and execution of marine operations. The importance of such work cannot be over-emphasised: In tonnage terms, 95 per cent of all international transport is by sea, and 20 per cent of the world’s oil and gas is produced from subsea reservoirs via offshore structures and pipelines. In the future, food production in aquacultural plants and the exploitation of renewable energy from the oceans is expected to play a growing role. The scientific and engineering research carried out in the centre takes account of such future needs, and extends current knowledge in relevant disciplines. The emphasis is on the synergy between hydrodynamics, structural mechanics and automatic control.
Facts and Figures
Personnel at the end of 2012:
6 key professors,
2 administrative staff,
11 postdocs/researchers,
43 PhD candidates,
7 graduated PhD candidates,
28 graduated MSc students,
2 adjunct professors,
3 adjunct associate professors,
4 visiting professors (short term),
3 visiting INSEAN and MIT researchers (short term)
and 5 visiting PhD candidates.
Revenues in 2012:
Income NOK 43 388 000,
Costs NOK 43 198 000.
Publications in 2012:
2 books,
14 book chapters,
10 international keynote lectures,
57 refereed journal papers
and 95 refereed conference papers.
Commitment of the key persons to international journals, conferences and workshops in 2012:
21 editorial boards of journals
and 14 international conference organising committees.
References
External links
CeSOS på NTNU's nettsted
Engineering research institutes
Marine engineering organizations | Centre for Ships and Ocean Structures | [
"Engineering"
] | 511 | [
"Marine engineering organizations",
"Engineering research institutes",
"Marine engineering"
] |
36,224,143 | https://en.wikipedia.org/wiki/Planetary%20science | Planetary science (or more rarely, planetology) is the scientific study of planets (including Earth), celestial bodies (such as moons, asteroids, comets) and planetary systems (in particular those of the Solar System) and the processes of their formation. It studies objects ranging in size from micrometeoroids to gas giants, with the aim of determining their composition, dynamics, formation, interrelations and history. It is a strongly interdisciplinary field, which originally grew from astronomy and Earth science, and now incorporates many disciplines, including planetary geology, cosmochemistry, atmospheric science, physics, oceanography, hydrology, theoretical planetary science, glaciology, and exoplanetology. Allied disciplines include space physics, when concerned with the effects of the Sun on the bodies of the Solar System, and astrobiology.
There are interrelated observational and theoretical branches of planetary science. Observational research can involve combinations of space exploration, predominantly with robotic spacecraft missions using remote sensing, and comparative, experimental work in Earth-based laboratories. The theoretical component involves considerable computer simulation and mathematical modelling.
Planetary scientists are generally located in the astronomy and physics or Earth sciences departments of universities or research centres, though there are several purely planetary science institutes worldwide. Generally, planetary scientists study one of the Earth sciences, astronomy, astrophysics, geophysics, or physics at the graduate level and concentrate their research in planetary science disciplines. There are several major conferences each year, and a wide range of peer reviewed journals. Some planetary scientists work at private research centres and often initiate partnership research tasks.
History
The history of planetary science may be said to have begun with the Ancient Greek philosopher Democritus, who is reported by Hippolytus as saying The ordered worlds are boundless and differ in size, and that in some there is neither sun nor moon, but that in others, both are greater than with us, and yet with others more in number. And that the intervals between the ordered worlds are unequal, here more and there less, and that some increase, others flourish and others decay, and here they come into being and there they are eclipsed. But that they are destroyed by colliding with one another. And that some ordered worlds are bare of animals and plants and all water.
In more modern times, planetary science began in astronomy, from studies of the unresolved planets. In this sense, the original planetary astronomer would be Galileo, who discovered the four largest moons of Jupiter, the mountains on the Moon, and first observed the rings of Saturn, all objects of intense later study. Galileo's study of the lunar mountains in 1609 also began the study of extraterrestrial landscapes: his observation "that the Moon certainly does not possess a smooth and polished surface" suggested that it and other worlds might appear "just like the face of the Earth itself".
Advances in telescope construction and instrumental resolution gradually allowed increased identification of the atmospheric as well as surface details of the planets. The Moon was initially the most heavily studied, due to its proximity to the Earth, as it always exhibited elaborate features on its surface, and the technological improvements gradually produced more detailed lunar geological knowledge. In this scientific process, the main instruments were astronomical optical telescopes (and later radio telescopes) and finally robotic exploratory spacecraft, such as space probes.
The Solar System has now been relatively well-studied, and a good overall understanding of the formation and evolution of this planetary system exists. However, there are large numbers of unsolved questions, and the rate of new discoveries is very high, partly due to the large number of interplanetary spacecraft currently exploring the Solar System.
Disciplines
Planetary science studies observational and theoretical astronomy, geology (astrogeology), atmospheric science, and an emerging subspecialty in planetary oceans, called planetary oceanography.
Planetary astronomy
This is both an observational and a theoretical science. Observational researchers are predominantly concerned with the study of the small bodies of the Solar System: those that are observed by telescopes, both optical and radio, so that characteristics of these bodies such as shape, spin, surface materials and weathering are determined, and the history of their formation and evolution can be understood.
Theoretical planetary astronomy is concerned with dynamics: the application of the principles of celestial mechanics to the Solar System and extrasolar planetary systems. Observing exoplanets and determining their physical properties, exoplanetology, is a major area of research besides Solar System studies. Every planet has its own branch.
Planetary geology
In planetary science, the term geology is used in its broadest sense, to mean the study of the surface and interior parts of planets and moons, from their core to their magnetosphere. The best-known research topics of planetary geology deal with the planetary bodies in the near vicinity of the Earth: the Moon, and the two neighboring planets: Venus and Mars. Of these, the Moon was studied first, using methods developed earlier on the Earth. Planetary geology focuses on celestial objects that exhibit a solid surface or have significant solid physical states as part of their structure. Planetary geology applies geology, geophysics and geochemistry to planetary bodies.
Planetary geomorphology
Geomorphology studies the features on planetary surfaces and reconstructs the history of their formation, inferring the physical processes that acted on the surface. Planetary geomorphology includes the study of several classes of surface features:
Impact features (multi-ringed basins, craters)
Volcanic and tectonic features (lava flows, fissures, rilles)
Glacial features
Aeolian features
Space weathering – erosional effects generated by the harsh environment of space (continuous micrometeorite bombardment, high-energy particle rain, impact gardening). For example, the thin dust cover on the surface of the lunar regolith is a result of micrometeorite bombardment.
Hydrological features: the liquid involved can range from water to hydrocarbon and ammonia, depending on the location within the Solar System. This category includes the study of paleohydrological features (paleochannels, paleolakes).
The history of a planetary surface can be deciphered by mapping features from top to bottom according to their deposition sequence, as first determined on terrestrial strata by Nicolas Steno. For example, stratigraphic mapping prepared the Apollo astronauts for the field geology they would encounter on their lunar missions. Overlapping sequences were identified on images taken by the Lunar Orbiter program, and these were used to prepare a lunar stratigraphic column and geological map of the Moon.
Cosmochemistry, geochemistry and petrology
One of the main problems when generating hypotheses on the formation and evolution of objects in the Solar System is the lack of samples that can be analyzed in the laboratory, where a large suite of tools are available, and the full body of knowledge derived from terrestrial geology can be brought to bear. Direct samples from the Moon, asteroids and Mars are present on Earth, removed from their parent bodies, and delivered as meteorites. Some of these have suffered contamination from the oxidising effect of Earth's atmosphere and the infiltration of the biosphere, but those meteorites collected in the last few decades from Antarctica are almost entirely pristine.
The different types of meteorites that originate from the asteroid belt cover almost all parts of the structure of differentiated bodies: meteorites even exist that come from the core-mantle boundary (pallasites). The combination of geochemistry and observational astronomy has also made it possible to trace the HED meteorites back to a specific asteroid in the main belt, 4 Vesta.
The comparatively few known Martian meteorites have provided insight into the geochemical composition of the Martian crust, although the unavoidable lack of information about their points of origin on the diverse Martian surface has meant that they do not provide more detailed constraints on theories of the evolution of the Martian lithosphere. As of July 24, 2013, 65 samples of Martian meteorites have been discovered on Earth. Many were found in either Antarctica or the Sahara Desert.
During the Apollo era, in the Apollo program, 384 kilograms of lunar samples were collected and transported to the Earth, and three Soviet Luna robots also delivered regolith samples from the Moon. These samples provide the most comprehensive record of the composition of any Solar System body besides the Earth. The numbers of lunar meteorites are growing quickly in the last few years – as of
April 2008 there are 54 meteorites that have been officially classified as lunar.
Eleven of these are from the US Antarctic meteorite collection, 6 are from the Japanese
Antarctic meteorite collection and the other 37 are from hot desert localities in Africa,
Australia, and the Middle East. The total mass of recognized lunar meteorites is close to
50 kg.
Planetary geophysics and space physics
Space probes made it possible to collect data in not only the visible light region but in other areas of the electromagnetic spectrum. The planets can be characterized by their force fields: gravity and their magnetic fields, which are studied through geophysics and space physics.
Measuring the changes in acceleration experienced by spacecraft as they orbit has allowed fine details of the gravity fields of the planets to be mapped. For example, in the 1970s, the gravity field disturbances above lunar maria were measured through lunar orbiters, which led to the discovery of concentrations of mass, mascons, beneath the Imbrium, Serenitatis, Crisium, Nectaris and Humorum basins.
If a planet's magnetic field is sufficiently strong, its interaction with the solar wind forms a magnetosphere around a planet. Early space probes discovered the gross dimensions of the terrestrial magnetic field, which extends about 10 Earth radii towards the Sun. The solar wind, a stream of charged particles, streams out and around the terrestrial magnetic field, and continues behind the magnetic tail, hundreds of Earth radii downstream. Inside the magnetosphere, there are relatively dense regions of solar wind particles, the Van Allen radiation belts.
Planetary geophysics includes, but is not limited to, seismology and tectonophysics, geophysical fluid dynamics, mineral physics, geodynamics, mathematical geophysics, and geophysical surveying.
Planetary geodesy
Planetary geodesy (also known as planetary geodetics) deals with the measurement and representation of the planets of the Solar System, their gravitational fields and geodynamic phenomena (polar motion in three-dimensional, time-varying space). The science of geodesy has elements of both astrophysics and planetary sciences. The shape of the Earth is to a large extent the result of its rotation, which causes its equatorial bulge, and the competition of geologic processes such as the collision of plates and of vulcanism, resisted by the Earth's gravity field. These principles can be applied to the solid surface of Earth (orogeny; Few mountains are higher than , few deep sea trenches deeper than that because quite simply, a mountain as tall as, for example, , would develop so much pressure at its base, due to gravity, that the rock there would become plastic, and the mountain would slump back to a height of roughly in a geologically insignificant time. Some or all of these geologic principles can be applied to other planets besides Earth. For instance on Mars, whose surface gravity is much less, the largest volcano, Olympus Mons, is high at its peak, a height that could not be maintained on Earth. The Earth geoid is essentially the figure of the Earth abstracted from its topographic features. Therefore, the Mars geoid (areoid) is essentially the figure of Mars abstracted from its topographic features. Surveying and mapping are two important fields of application of geodesy.
Planetary atmospheric science
An atmosphere is an important transitional zone between the solid planetary surface and the higher rarefied ionizing and radiation belts. Not all planets have atmospheres: their existence depends on the mass of the planet, and the planet's distance from the Sun – too distant and frozen atmospheres occur. Besides the four giant planets, three of the four terrestrial planets (Earth, Venus, and Mars) have significant atmospheres. Two moons have significant atmospheres: Saturn's moon Titan and Neptune's moon Triton. A tenuous atmosphere exists around Mercury.
The effects of the rotation rate of a planet about its axis can be seen in atmospheric streams and currents. Seen from space, these features show as bands and eddies in the cloud system and are particularly visible on Jupiter and Saturn.
Planetary oceanography
Exoplanetology
Exoplanetology studies exoplanets, the planets existing outside our Solar System. Until recently, the means of studying exoplanets have been extremely limited, but with the current rate of innovation in research technology, exoplanetology has become a rapidly developing subfield of astronomy.
Comparative planetary science
Planetary science frequently makes use of the method of comparison to give a greater understanding of the object of study. This can involve comparing the dense atmospheres of Earth and Saturn's moon Titan, the evolution of outer Solar System objects at different distances from the Sun, or the geomorphology of the surfaces of the terrestrial planets, to give only a few examples.
The main comparison that can be made is to features on the Earth, as it is much more accessible and allows a much greater range of measurements to be made. Earth analog studies are particularly common in planetary geology, geomorphology, and also in atmospheric science.
The use of terrestrial analogs was first described by Gilbert (1886).
In fiction
In Frank Herbert's 1965 science fiction novel Dune, the major secondary character Liet-Kynes serves as the "Imperial Planetologist" for the fictional planet Arrakis, a position he inherited from his father Pardot Kynes. In this role, a planetologist is described as having skills of an ecologist, geologist, meteorologist, and biologist, as well as basic understandings of human sociology. The planetologists apply this expertise to the study of entire planets.In the Dune series, planetologists are employed to understand planetary resources and to plan terraforming or other planetary-scale engineering projects. This fictional position in Dune has had an impact on the discourse surrounding planetary science itself and is referred to by one author as a "touchstone" within the related disciplines. In one example, a publication by Sybil P. Seitzinger in the journal Nature opens with a brief introduction on the fictional role in Dune, and suggests we should consider appointing individuals with similar skills to Liet-Kynes to help with managing human activity on Earth.
Professional activity
Journals
Annual Review of Earth and Planetary Sciences
Earth and Planetary Science Letters
Earth, Moon, and Planets
Geochimica et Cosmochimica Acta
Icarus
Journal of Geophysical Research – Planets
Meteoritics and Planetary Science
Planetary and Space Science
The Planetary Science Journal
Professional bodies
This non-exhaustive list includes those institutions and universities with major groups of people working in planetary science. Alphabetical order is used.
Division for Planetary Sciences (DPS) of the American Astronomical Society
American Geophysical Union
Meteoritical Society
Europlanet
Government space agencies
Canadian Space Agency (CSA)
China National Space Administration (CNSA, People's Republic of China).
French National Centre of Space Research
Deutsches Zentrum für Luft- und Raumfahrt e.V., (German: abbreviated DLR), the German Aerospace Center
European Space Agency (ESA)
Indian Space Research Organisation (ISRO)
Israel Space Agency (ISA)
Italian Space Agency
Japan Aerospace Exploration Agency (JAXA)
NASA (National Aeronautics and Space Administration, United States of America)
JPL
GSFC
Ames
National Space Organization (Taiwan).
Russian Federal Space Agency
UK Space Agency (UKSA).
Major conferences
Lunar and Planetary Science Conference (LPSC), organized by the Lunar and Planetary Institute in Houston. Held annually since 1970, occurs in March.
Division for Planetary Sciences (DPS) meeting held annually since 1970 at a different location each year, predominantly within the mainland US. Occurs around October.
American Geophysical Union (AGU) annual Fall meeting in December in San Francisco.
American Geophysical Union (AGU) Joint Assembly (co-sponsored with other societies) in April–May, in various locations around the world.
Meteoritical Society annual meeting, held during the Northern Hemisphere summer, generally alternating between North America and Europe.
European Planetary Science Congress (EPSC), held annually around September at a location within Europe.
Smaller workshops and conferences on particular fields occur worldwide throughout the year.
See also
Areography (geography of Mars)
Planetary cartography
Planetary coordinate system
Selenography – study of the surface and physical features of the Moon
Theoretical planetology
Timeline of Solar System exploration
References
Further reading
Carr, Michael H., Saunders, R. S., Strom, R. G., Wilhelms, D. E. 1984. The Geology of the Terrestrial Planets. NASA.
Morrison, David. 1994. Exploring Planetary Worlds. W. H. Freeman.
Hargitai H et al. (2015) Classification and Characterization of Planetary Landforms. In: Hargitai H (ed) Encyclopedia of Planetary Landforms. Springer. https://link.springer.com/content/pdf/bbm%3A978-1-4614-3134-3%2F1.pdf
Hauber E et al. (2019) Planetary geologic mapping. In: Hargitai H (ed) Planetary Cartography and GIS. Springer.
Page D (2015) The Geology of Planetary Landforms. In: Hargitai H (ed) Encyclopedia of Planetary Landforms. Springer.
Rossi, A.P., van Gasselt S (eds) (2018) Planetary Geology. Springer
External links
Planetary Science Research Discoveries (articles)
The Planetary Society (world's largest space-interest group: see also their active news blog)
Planetary Exploration Newsletter (PSI-published professional newsletter, weekly distribution)
Women in Planetary Science (professional networking and news)
Space science
Astronomical sub-disciplines | Planetary science | [
"Astronomy"
] | 3,678 | [
"Space science",
"Planetary science",
"Outer space",
"Astronomical sub-disciplines"
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36,225,667 | https://en.wikipedia.org/wiki/Schwinger%20limit | In quantum electrodynamics (QED), the Schwinger limit is a scale above which the electromagnetic field is expected to become nonlinear. The limit was first derived in one of QED's earliest theoretical successes by Fritz Sauter in 1931 and discussed further by Werner Heisenberg and his student Hans Heinrich Euler. The limit, however, is commonly named in the literature for Julian Schwinger, who derived the leading nonlinear corrections to the fields and calculated the rate of electron–positron pair production in a strong electric field. The limit is typically reported as a maximum electric field or magnetic field before nonlinearity for the vacuum of
where me is the mass of the electron, c is the speed of light in vacuum, qe is the elementary charge, and ħ is the reduced Planck constant. These are enormous field strengths. Such an electric field is capable of accelerating a proton from rest to the maximum energy attained by protons at the Large Hadron Collider in only approximately 5 micrometers. The magnetic field is associated with birefringence of the vacuum and is exceeded on magnetars.
In vacuum, the classical Maxwell's equations are perfectly linear differential equations. This implies – by the superposition principle – that the sum of any two solutions to Maxwell's equations is another solution to Maxwell's equations. For example, two intersecting beams of light should simply add together their electric fields and pass right through each other. Thus Maxwell's equations predict the impossibility of any but trivial elastic photon–photon scattering. In QED, however, non-elastic photon–photon scattering becomes possible when the combined energy is large enough to create virtual electron–positron pairs spontaneously, illustrated by the Feynman diagram in the adjacent figure. This creates nonlinear effects that are approximately described by Euler and Heisenberg's nonlinear variant of Maxwell's equations.
A single plane wave is insufficient to cause nonlinear effects, even in QED. The basic reason for this is that a single plane wave of a given energy may always be viewed in a different reference frame, where it has less energy (the same is the case for a single photon). A single wave or photon does not have a center-of-momentum frame where its energy must be at minimal value. However, two waves or two photons not traveling in the same direction always have a minimum combined energy in their center-of-momentum frame, and it is this energy and the electric field strengths associated with it, which determine particle–antiparticle creation, and associated scattering phenomena.
Photon–photon scattering and other effects of nonlinear optics in vacuum is an active area of experimental research, with current or planned technology beginning to approach the Schwinger limit. It has already been observed through inelastic channels in SLAC Experiment 144. However, the direct effects in elastic scattering have not been observed. As of 2012, the best constraint on the elastic photon–photon scattering cross section belonged to PVLAS, which reported an upper limit far above the level predicted by the Standard Model.
Proposals were made to measure elastic light-by-light scattering using the strong electromagnetic fields of the hadrons collided at the LHC. In 2019, the ATLAS experiment at the LHC announced the first definitive observation of photon–photon scattering, observed in lead ion collisions that produced fields as large as , well in excess of the Schwinger limit. Observation of a cross section larger or smaller than that predicted by the Standard Model could signify new physics such as axions, the search of which is the primary goal of PVLAS and several similar experiments. ATLAS observed more events than expected, potentially evidence that the cross section is larger than predicted by the Standard Model, but the excess is not yet statistically significant.
The planned, funded ELI–Ultra High Field Facility, which will study light at the intensity frontier, is likely to remain well below the Schwinger limit although it may still be possible to observe some nonlinear optical effects. The Station of Extreme Light (SEL) is another laser facility under construction which should be powerful enough to observe the effect. Such an experiment, in which ultra-intense light causes pair production, has been described in the popular media as creating a "hernia" in spacetime.
See also
Julian Schwinger
Schwinger effect
Sokolov–Ternov effect
Vacuum polarization
References
Particle physics
Quantum electrodynamics
Quantum optics | Schwinger limit | [
"Physics"
] | 906 | [
"Quantum optics",
"Quantum mechanics",
"Particle physics"
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36,227,078 | https://en.wikipedia.org/wiki/Oceanic%20carbon%20cycle | The oceanic carbon cycle (or marine carbon cycle) is composed of processes that exchange carbon between various pools within the ocean as well as between the atmosphere, Earth interior, and the seafloor. The carbon cycle is a result of many interacting forces across multiple time and space scales that circulates carbon around the planet, ensuring that carbon is available globally. The Oceanic carbon cycle is a central process to the global carbon cycle and contains both inorganic carbon (carbon not associated with a living thing, such as carbon dioxide) and organic carbon (carbon that is, or has been, incorporated into a living thing). Part of the marine carbon cycle transforms carbon between non-living and living matter.
Three main processes (or pumps) that make up the marine carbon cycle bring atmospheric carbon dioxide (CO2) into the ocean interior and distribute it through the oceans. These three pumps are: (1) the solubility pump, (2) the carbonate pump, and (3) the biological pump. The total active pool of carbon at the Earth's surface for durations of less than 10,000 years is roughly 40,000 gigatons C (Gt C, a gigaton is one billion tons, or the weight of approximately 6 million blue whales), and about 95% (~38,000 Gt C) is stored in the ocean, mostly as dissolved inorganic carbon. The speciation of dissolved inorganic carbon in the marine carbon cycle is a primary controller of acid-base chemistry in the oceans.
Earth's plants and algae (primary producers) are responsible for the largest annual carbon fluxes. Although the amount of carbon stored in marine biota (~3 Gt C) is very small compared with terrestrial vegetation (~610 GtC), the amount of carbon exchanged (the flux) by these groups is nearly equal – about 50 GtC each. Marine organisms link the carbon and oxygen cycles through processes such as photosynthesis. The marine carbon cycle is also biologically tied to the nitrogen and phosphorus cycles by a near-constant stoichiometric ratio C:N:P of 106:16:1, also known as the Redfield Ketchum Richards (RKR) ratio, which states that organisms tend to take up nitrogen and phosphorus incorporating new organic carbon. Likewise, organic matter decomposed by bacteria releases phosphorus and nitrogen.
Based on the publications of NASA, World Meteorological Association, IPCC, and International Council for the Exploration of the Sea, as well as scientists from NOAA, Woods Hole Oceanographic Institution, Scripps Institution of Oceanography, CSIRO, and Oak Ridge National Laboratory, the human impacts on the marine carbon cycle are significant. Before the Industrial Revolution, the ocean was a net source of CO2 to the atmosphere whereas now the majority of the carbon that enters the ocean comes from atmospheric carbon dioxide (CO2).
In recent decades, the ocean has acted as a sink for anthropogenic CO2, absorbing around a quarter of the CO2 produced by humans through the burning of fossil fuels and land use changes. By doing so, the ocean has acted as a buffer, somewhat slowing the rise in atmospheric CO2 levels. However, this absorption of anthropogenic CO2 has also caused acidification of the oceans. Climate change, a result of this excess CO2 in the atmosphere, has increased the temperature of the ocean and atmosphere. The slowed rate of global warming occurring from 2000–2010 may be attributed to an observed increase in upper ocean heat content.
Marine carbon
Carbon compounds can be distinguished as either organic or inorganic, and dissolved or particulate, depending on their composition. Organic carbon forms the backbone of key component of organic compounds such as – proteins, lipids, carbohydrates, and nucleic acids. Inorganic carbon is found primarily in simple compounds such as carbon dioxide, carbonic acid, bicarbonate, and carbonate (CO2, H2CO3, HCO3−, CO32− respectively).
Marine carbon is further separated into particulate and dissolved phases. These pools are operationally defined by physical separation – dissolved carbon passes through a 0.2 μm filter, and particulate carbon does not.
Inorganic carbon
There are two main types of inorganic carbon that are found in the oceans. Dissolved inorganic carbon (DIC) is made up of bicarbonate (HCO3−), carbonate (CO32−) and carbon dioxide (including both dissolved CO2 and carbonic acid H2CO3). DIC can be converted to particulate inorganic carbon (PIC) through precipitation of CaCO3 (biologically or abiotically). DIC can also be converted to particulate organic carbon (POC) through photosynthesis and chemoautotrophy (i.e. primary production). DIC increases with depth as organic carbon particles sink and are respired. Free oxygen decreases as DIC increases because oxygen is consumed during aerobic respiration.
Particulate inorganic carbon (PIC) is the other form of inorganic carbon found in the ocean. Most PIC is the CaCO3 that makes up shells of various marine organisms, but can also form in whiting events. Marine fish also excrete calcium carbonate during osmoregulation.
Some of the inorganic carbon species in the ocean, such as bicarbonate and carbonate, are major contributors to alkalinity, a natural ocean buffer that prevents drastic changes in acidity (or pH). The marine carbon cycle also affects the reaction and dissolution rates of some chemical compounds, regulates the amount of carbon dioxide in the atmosphere and Earth's temperature.
Organic carbon
Like inorganic carbon, there are two main forms of organic carbon found in the ocean (dissolved and particulate). Dissolved organic carbon (DOC) is defined operationally as any organic molecule that can pass through a 0.2 μm filter. DOC can be converted into particulate organic carbon through heterotrophy and it can also be converted back to dissolved inorganic carbon (DIC) through respiration.
Those organic carbon molecules being captured on a filter are defined as particulate organic carbon (POC). POC is composed of organisms (dead or alive), their fecal matter, and detritus. POC can be converted to DOC through disaggregation of molecules and by exudation by phytoplankton, for example. POC is generally converted to DIC through heterotrophy and respiration.
Marine carbon pumps
Solubility pump
Full article: Solubility pump
The oceans store the largest pool of reactive carbon on the planet as DIC, which is introduced as a result of the dissolution of atmospheric carbon dioxide into seawater – the solubility pump. Aqueous CO2, carbonic acid, bicarbonate ion, and carbonate ion concentrations comprise dissolved inorganic carbon (DIC). DIC circulates throughout the whole ocean by Thermohaline circulation, which facilitates the tremendous DIC storage capacity of the ocean. The chemical equations below show the reactions that CO2 undergoes after it enters the ocean and transforms into its aqueous form.
Carbonic acid rapidly dissociates into free hydrogen ion (technically, hydronium) and bicarbonate.The free hydrogen ion meets carbonate, already present in the water from the dissolution of CaCO3, and reacts to form more bicarbonate ion.The dissolved species in the equations above, mostly bicarbonate, make up the carbonate alkalinity system, the dominant contributor to seawater alkalinity.
Carbonate pump
The carbonate pump, sometimes called the carbonate counter pump, starts with marine organisms at the ocean's surface producing particulate inorganic carbon (PIC) in the form of calcium carbonate (calcite or aragonite, CaCO3). This CaCO3 is what forms hard body parts like shells. The formation of these shells increases atmospheric CO2 due to the production of CaCO3 in the following reaction with simplified stoichiometry:Coccolithophores, a nearly ubiquitous group of phytoplankton that produce shells of calcium carbonate, are the dominant contributors to the carbonate pump. Due to their abundance, coccolithophores have significant implications on carbonate chemistry, in the surface waters they inhabit and in the ocean below: they provide a large mechanism for the downward transport of CaCO3. The air-sea CO2 flux induced by a marine biological community can be determined by the rain ratio - the proportion of carbon from calcium carbonate compared to that from organic carbon in particulate matter sinking to the ocean floor, (PIC/POC). The carbonate pump acts as a negative feedback on CO2 taken into the ocean by the solubility pump. It occurs with lesser magnitude than the solubility pump.
Biological pump
Full article: Biological pump
Particulate organic carbon, created through biological production, can be exported from the upper ocean in a flux commonly termed the biological pump, or respired (equation 6) back into inorganic carbon. In the former, dissolved inorganic carbon is biologically converted into organic matter by photosynthesis (equation 5) and other forms of autotrophy that then sinks and is, in part or whole, digested by heterotrophs. Particulate organic carbon can be classified, based on how easily organisms can break them down for food, as labile, semilabile, or refractory. Photosynthesis by phytoplankton is the primary source for labile and semilabile molecules, and is the indirect source for most refractory molecules. Labile molecules are present at low concentrations outside of cells (in the picomolar range) and have half-lives of only minutes when free in the ocean. They are consumed by microbes within hours or days of production and reside in the surface oceans, where they contribute a majority of the labile carbon flux. Semilabile molecules, much more difficult to consume, are able to reach depths of hundreds of meters below the surface before being metabolized. Refractory DOM largely comprises highly conjugated molecules like Polycyclic aromatic hydrocarbons or lignin. Refractory DOM can reach depths greater than 1000 m and circulates through the oceans over thousands of years. Over the course of a year, approximately 20 gigatons of photosynthetically-fixed labile and semilabile carbon is taken up by heterotrophs, whereas fewer than 0.2 gigatons of refractory carbon is consumed. Marine dissolved organic matter (DOM) can store as much carbon as the current atmospheric CO2 supply, but industrial processes are altering the balance of this cycle.
Inputs
Inputs to the marine carbon cycle are numerous, but the primary contributions, on a net basis, come from the atmosphere and rivers. Hydrothermal vents generally supply carbon equal to the amount they consume.
Atmosphere
Before the Industrial Revolution, the ocean was a source of CO2 to the atmosphere balancing the impact of rock weathering and terrestrial particulate organic carbon; now it has become a sink for the excess atmospheric CO2. Carbon dioxide is absorbed from the atmosphere at the ocean's surface at an exchange rate which varies locally and with time but on average, the oceans have a net absorption of around 2.9 Pg (equivalent to 2.9 billion metric tonnes) of carbon from atmospheric CO2 per year. Because the solubility of carbon dioxide increases when temperature decreases, cold areas can contain more CO2 and still be in equilibrium with the atmosphere; In contrast, rising sea surface temperatures decrease the capacity of the oceans to take in carbon dioxide. The North Atlantic and Nordic oceans have the highest carbon uptake per unit area in the world, and in the North Atlantic deep convection transports approximately 197 Tg per year of non-refractory carbon to depth.
The rate of CO2 absorption by the ocean has been increasing with time as atmospheric CO2 concentrations have increased due to anthropogenic emissions. However, the ocean carbon sink may be more sensitive to climate change than previously thought, and ocean warming and circulation changes due to climate change could result in the ocean absorbing less CO2 from the atmosphere in future than expected.
Carbon dioxide exchange rates between ocean and atmosphere
Ocean-atmospheric exchanges rates of CO2 depend on the concentration of carbon dioxide already present in both the atmosphere and the ocean, temperature, salinity, and wind speed. This exchange rate can be approximated by Henry's law and can be calculated as S = kP, where the solubility (S) of the carbon dioxide gas is proportional to the amount of gas in the atmosphere, or its partial pressure.
Revelle factor
Since the oceanic intake of carbon dioxide is limited, CO2 influx can also be described by the Revelle factor. The Revelle Factor is a ratio of the change of carbon dioxide to the change in dissolved inorganic carbon, which serves as an indicator of carbon dioxide dissolution in the mixed layer considering the solubility pump. The Revelle Factor is an expression to characterize the thermodynamic efficiency of the DIC pool to absorb CO2 into bicarbonate. The lower the Revelle factor, the higher the capacity for ocean water to take in carbon dioxide. While Revelle calculated a factor of around 10 in his day, in a 2004 study data showed a Revelle factor ranging from approximately 9 in low-latitude tropical regions to 15 in the southern ocean near Antarctica.
Rivers
Rivers can also transport organic carbon to the ocean through weathering or erosion of aluminosilicate (equation 7) and carbonate rocks (equation 8) on land,
or by the decomposition of life (equation 5, e.g. plant and soil material). Rivers contribute roughly equal amounts (~0.4 GtC/yr) of DIC and DOC to the oceans. It is estimated that approximately 0.8 GtC (DIC + DOC) is transported annually from the rivers to the ocean. The rivers that flow into Chesapeake Bay (Susquehanna, Potomac, and James rivers) input approximately 0.004 Gt (6.5 x 1010 moles) DIC per year. The total carbon transport of rivers represents approximately 0.02% of the total carbon in the atmosphere. Though it seems small, over long time scales (1000 to 10,000 years) the carbon that enters rivers (and therefore does not enter the atmosphere) serves as a stabilizing feedback for greenhouse warming.
Outputs
The key outputs of the marine carbon system are particulate organic matter (POC) and calcium carbonate (PIC) preservation as well as reverse weathering. While there are regions with local loss of CO2 to the atmosphere and hydrothermal processes, a net loss in the cycle does not occur.
Organic matter preservation
Sedimentation is a long-term sink for carbon in the ocean, as well as the largest loss of carbon from the oceanic system. Deep marine sediments and geologic formations are important since they provide a thorough record of life on Earth and an important source of fossil fuel. Oceanic carbon can exit the system in the form of detritus that sinks and is buried in the seafloor without being fully decomposed or dissolved. Ocean floor surface sediments account for 1.75x1015 kg of carbon in the global carbon cycle. At most, 4% of the particulate organic carbon from the euphotic zone in the Pacific Ocean, where light-powered primary production occurs, is buried in marine sediments. It is then implied that since there is a higher input of organic matter to the ocean than what is being buried, a large portion of it is used up or consumed within.
Fate of sinking organic carbon
Historically, sediments with the highest organic carbon contents were frequently found in areas with high surface water productivity or those with low bottom-water oxygen concentrations. 90% of organic carbon burial occurs in deposits of deltas and continental shelves and upper slopes; this is due partly to short exposure time because of a shorter distance to the seafloor and the composition of the organic matter that is already deposited in those environments. Organic carbon burial is also sensitive to climate patterns: the accumulation rate of organic carbon was 50% larger during the glacial maximum compared to interglacials.
Degradation
POC is decomposed by a series of microbe-driven processes, such as methanogenesis and sulfate reduction, before burial in the seafloor. Degradation of POC also results in microbial methane production which is the main gas hydrate on the continental margins. Lignin and pollen are inherently resistant to degradation, and some studies show that inorganic matrices may also protect organic matter. Preservation rates of organic matter depend on other interdependent variables that vary nonlinearly in time and space. Although organic matter breakdown occurs rapidly in the presence of oxygen, microbes utilizing a variety of chemical species (via redox gradients) can degrade organic matter in anoxic sediments. The burial depth at which degradation halts depends upon the sedimentation rate, the relative abundance of organic matter in the sediment, the type of organic matter being buried, and innumerable other variables. While decomposition of organic matter can occur in anoxic sediments when bacteria use oxidants other than oxygen (nitrate, sulfate, Fe3+), decomposition tends to end short of complete mineralization. This occurs because of preferential decomposition of labile molecules over refractile molecules.
Burial
Organic carbon burial is an input of energy for underground biological environments and can regulate oxygen in the atmosphere at long time-scales (> 10,000 years). Burial can only take place if organic carbon arrives to the sea floor, making continental shelves and coastal margins the main storage of organic carbon from terrestrial and oceanic primary production. Fjords, or cliffs created by glacial erosion, have also been identified as areas of significant carbon burial, with rates one hundred times greater than the ocean average. Particulate organic carbon is buried in oceanic sediments, creating a pathway between a rapidly available carbon pool in the ocean to its storage for geological timescales. Once carbon is sequestered in the seafloor, it is considered blue carbon. Burial rates can be calculated as the difference between the rate at which organic matter sinks and the rate at which it decomposes.
Calcium carbonate preservation
The precipitation of calcium carbonate is important as it results in a loss of alkalinity as well as a release of CO2 (Equation 4), and therefore a change in the rate of preservation of calcium carbonate can alter the partial pressure of CO2 in Earth's atmosphere. CaCO3 is supersatured in the great majority of ocean surface waters and undersaturated at depth, meaning the shells are more likely to dissolve as they sink to ocean depths. CaCO3 can also be dissolved through metabolic dissolution (i.e. can be used as food and excreted) and thus deep ocean sediments have very little calcium carbonate. The precipitation and burial of calcium carbonate in the ocean removes particulate inorganic carbon from the ocean and ultimately forms limestone. On time scales greater than 500,000 years Earth's climate is moderated by the flux of carbon in and out of the lithosphere. Rocks formed in the ocean seafloor are recycled through plate tectonics back to the surface and weathered or subducted into the mantle, the carbon outgassed by volcanoes.
Human impacts
Oceans take up around 25 – 31% of anthropogenic CO2. Because the Revelle factor increases with increasing CO2, a smaller fraction of the anthropogenic flux will be taken up by the ocean in the future. Current annual increase in atmospheric CO2 is approximately 4–5 gigatons of carbon, about 2–3ppm CO2 per year. This induces climate change that drives carbon concentration and carbon-climate feedback processes that modifies ocean circulation and the physical and chemical properties of seawater, which alters CO2 uptake. Overfishing and the plastic pollution of the oceans contribute to the degraded state of the world's biggest carbon sink.
Ocean acidification
Full article: Ocean acidification
The pH of the oceans is declining due to uptake of atmospheric CO2. The rise in dissolved carbon dioxide reduces the availability of the carbonate ion, reducing CaCO3 saturation state, thus making it thermodynamically harder to make CaCO3 shell. Carbonate ions preferentially bind to hydrogen ions to form bicarbonate, thus a reduction in carbonate ion availability increases the amount of unbound hydrogen ions, and decreases the amount of bicarbonate formed (Equations 1–3). pH is a measurement of hydrogen ion concentration, where a low pH means there are more unbound hydrogen ions. pH is therefore an indicator of carbonate speciation (the format of carbon present) in the oceans and can be used to assess how healthy the ocean is.
The list of organisms that may struggle due to ocean acidification include coccolithophores and foraminifera (the base of the marine food chain in many areas), human food sources such as oysters and mussels, and perhaps the most conspicuous, a structure built by organisms – the coral reefs. Most surface water will remain supersaturated with respect to CaCO3 (both calcite and aragonite) for some time on current emissions trajectories, but the organisms that require carbonate will likely be replaced in many areas. Coral reefs are under pressure from overfishing, nitrate pollution, and warming waters; ocean acidification will add additional stress on these important structures.
Iron fertilization
Full article: Iron Fertilization
Iron fertilization is a facet of geoengineering, which purposefully manipulates the Earth's climate system, typically in aspects of the carbon cycle or radiative forcing. Of current geoengineering interest is the possibility of accelerating the biological pump to increase export of carbon from the surface ocean. This increased export could theoretically remove excess carbon dioxide from the atmosphere for storage in the deep ocean. Ongoing investigations regarding artificial fertilization exist. Due to the scale of the ocean and the fast response times of heterotrophic communities to increases in primary production, it is difficult to determine whether limiting-nutrient fertilization results in an increase in carbon export. However, the majority of the community does not believe this is a reasonable or viable approach.
Dams and reservoirs
There are over 16 million dams in the world that alter carbon transport from rivers to oceans. Using data from the Global Reservoirs and Dams database, which contains approximately 7000 reservoirs that hold 77% of the total volume of water held back by dams (8000 km3), it is estimated that the delivery of carbon to the ocean has decreased by 13% since 1970 and is projected to reach 19% by 2030. The excess carbon contained in the reservoirs may emit an additional ~0.184 Gt of carbon to the atmosphere per year and an additional ~0.2 GtC will be buried in sediment. Prior to 2000, the Mississippi, the Niger, and the Ganges River basins account for 25 – 31% of all reservoir carbon burial. After 2000, the Paraná (home to 70 dams) and the Zambezi (home to the largest reservoir) River basins exceeded the burial by the Mississippi. Other large contributors to carbon burial caused by damming occur on the Danube, the Amazon, the Yangtze, the Mekong, the Yenisei, and the Tocantins Rivers.
See also
Phosphorus cycle
References
External links
Current global map of the partial pressure of carbon dioxide at the ocean surface
Current global map of the sea-air carbon dioxide flux density
Carbon cycle
Chemical oceanography | Oceanic carbon cycle | [
"Chemistry"
] | 4,807 | [
"Chemical oceanography"
] |
36,229,076 | https://en.wikipedia.org/wiki/Multipass%20spectroscopic%20absorption%20cells | Multiple-pass or long path absorption cells are commonly used in spectroscopy to measure low-concentration components or to observe weak spectra in gases or liquids. Several important advances were made in this area beginning in the 1930s, and research into a wide range of applications continues to the present day.
Functional Overview
Generally the goal of this type of sample cell is to improve detection sensitivity by increasing the total optical path length that travels through a small, constant sample volume. In principle, a longer path length results in greater detection sensitivity. Focusing mirrors must be used to redirect the beam at each reflection point, resulting in the beam being restricted to a predefined space along a controlled path until it exits the optical cavity. The output of the cell is the input of an optical detector (a specialized type of transducer), which senses specific changes in the properties of the beam that occur during interaction with the test sample. For instance, the sample may absorb energy from the beam, resulting in an attenuation of the output that is detectable by the transducer. Two conventional multipass cells are called the White cell and Herriott cell.
Pfund Cell
In the late 1930s August Pfund used a triple-pass cell like the one shown above for atmospheric study. The cell, which became known as the Pfund cell, is constructed using two identical spherical mirrors, each having a hole carefully machined into its center. The separation distance between the mirrors is equal to the mirror focal length. A source enters from a hole in either mirror, is redirected twice at two reflection points, and then exits the cell through the other mirror on the third pass. The Pfund cell was one of the earliest examples of this type of spectroscopic technique and is noted for having used multiple passes.
White cell
The White cell was first described in 1942 by John U. White in his paper Long Optical Paths of Large Aperture, and was a significant improvement over previous long path spectroscopic measurement techniques. A White cell is constructed using three spherical, concave mirrors having the same radius of curvature. The mirrors are separated by a distance equal to their radii of curvature. The animation on the right shows a White Cell in which a beam makes eight reflective passes or traversals. The number of traversals can be changed quite easily by making slight rotational adjustments to either M2 or M3; however, the total number of traversals must always occur in multiples of four. The entering and exiting beams do not change position as traversals are added or removed, while the total number of traversals can be increased many times without changing the volume of the cell, and therefore the total optical path length can be made large compared to the volume of the sample under test. The spots from various passes can overlap on mirrors M2 and M3 but must be distinct on mirror M1. If the input beam is focused in the plane of M1, then each round trip will also be focused in this plane. The tighter the focus, the more nonoverlapping spots there can be on M1 and thus the higher the maximum pathlength.
At present the White cell is still the most commonly used multipass cell and provides many advantages. For example,
The number of traversals is easily controlled
It allows for high numerical aperture
It is reasonably stable (but not as stable as the Herriott cell)
White cells are available with path lengths ranging from less than a meter to many hundreds of meters.
Herriott cell
The Herriott cell first appeared in 1965 when Donald R. Herriott and Harry J. Schulte published Folded Optical Delay Lines while at Bell Laboratories. The Herriott cell is made up of two opposing spherical mirrors. A hole is machined into one of the mirrors to allow the input and output beams to enter and exit the cavity. Alternatively, the beam may exit through a hole in the opposite mirror. In this fashion the Herriott cell can support multiple light sources by providing multiple entrance and exit holes in either of the mirrors. Unlike the White cell, the number of traversals is controlled by adjusting the separation distance D between the two mirrors. This cell is also commonly used and has some advantages over the White cell:
It is simpler than the White cell with only two mirrors that are easier to position and less susceptible to mechanical disturbance of the cell
Can be more stable than the White cell
However, the Herriot cell does not accept high numerical aperture beams. In addition, larger sized mirrors must be used when longer path lengths are needed.
Circular Multipass Cells
Another category of multipass cells is generally referred to as circular multipass reflection cells. They were first introduced by Thoma and co-workers in 1994. Such cells rely on a circular arrangement of mirrors. The beam enters the cell under an angle and propagates on a star-shaped pattern (see picture on the right). The path length in circular multipass cells can be varied by adjusting the incidence angle of the beam. An advantage lies in their robustness towards mechanical stress such as vibrations or temperature changes. Furthermore, circular multipass cells stand out because of the small detection volumes they provide. A stable beam propagation is achieved by shaping individual reflection points to form a non-concentric mirror-arrangement.
In a special case, a circular mirror is used, allowing continuous adjustment of the angle of incidence. A drawback of this circular cell configuration is the inherent concentric mirror arrangement which leads to imperfect imaging after a large number of reflections.
See also
Laser absorption spectrometry
Tunable diode laser absorption spectroscopy
Optical system
Absorption spectroscopy
Infrared spectroscopy
Absorption (optics)
Optical density
Optical depth
Reflectivity
References
Optical devices
Physical chemistry
Spectroscopy | Multipass spectroscopic absorption cells | [
"Physics",
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26,370,510 | https://en.wikipedia.org/wiki/Levomilnacipran | Levomilnacipran (brand name Fetzima) is an antidepressant which was approved in the United States in 2013 for the treatment of major depressive disorder (MDD) in adults. It is the levorotatory enantiomer of milnacipran, and has similar effects and pharmacology, acting as a serotonin–norepinephrine reuptake inhibitor (SNRI).
Medical uses
Depression
The FDA approved levomilnacipran for treating major depressive disorder. This approval was based on the results of five clinical trials. The trials included one 10-week phase II and four 8-week phase III. Four of the five trials demonstrated a statistically significant superiority to placebo as measured by the Montgomery–Åsberg Depression Rating Scale. Superiority to placebo was also demonstrated by improvement in the Sheehan Disability Scale.
Side effects
Side effects seen more often with levomilnacipran than with placebo in clinical trials included nausea, dizziness, sweating, constipation, insomnia, increased heart rate and blood pressure, urinary hesitancy, erectile dysfunction and delayed ejaculation in males, vomiting, tachycardia, and palpitations.
Pharmacology
Pharmacodynamics
Relative to other SNRIs, levomilnacipran, as well as milnacipran, differ in that they are much more balanced reuptake inhibitors of serotonin and norepinephrine. To demonstrate, the serotonin:norepinephrine ratios of SNRIs are as follows: venlafaxine = 30:1, duloxetine = 10:1, desvenlafaxine = 14:1, milnacipran = 1.6:1, and levomilnacipran = 1:2. The clinical implications of more balanced elevations of serotonin and norepinephrine are unclear, but may include improved effectiveness, though also increased side effects.
Levomilnacipran is selective for the serotonin and norepinephrine transporters, lacking significant affinity for over 23 off-target sites. However, it does show some affinity for the dizocilpine (MK-801/) site of the NMDA receptor (Ki = 1.7 μM), and has been found to inhibit NR2A and NR2B subunit-containing NMDA receptors with respective IC50 values of 5.62 and 4.57 μM. As such, levomilnacipran is an NMDA receptor antagonist at high concentrations.
Levomilnacipran has recently been found to act as an inhibitor of beta-site amyloid precursor protein cleaving enzyme-1 (BACE-1), which is responsible for β-amyloid plaque formation, and hence may be a potentially useful drug in the treatment of Alzheimer's disease.
Pharmacokinetics
Levomilnacipran has a high oral bioavailability of 92% and a low plasma protein binding of 22%. It is metabolized in the liver by the cytochrome P450 enzyme CYP3A4, thereby making the medication susceptible to grapefruit-drug interactions. The drug has an elimination half-life of approximately 12 hours, allowing for once-daily administration. Levomilnacipran is excreted in urine.
History
Levomilnacipran was developed by Forest Laboratories and Pierre Fabre Group, and was approved by the Food and Drug Administration in July 2013.
References
External links
Alzheimer's disease
Carboxamides
Cyclopropanes
Drugs developed by AbbVie
Enantiopure drugs
Enzyme inhibitors
NMDA receptor antagonists
Phenethylamines
Serotonin–norepinephrine reuptake inhibitors | Levomilnacipran | [
"Chemistry"
] | 798 | [
"Stereochemistry",
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26,372,648 | https://en.wikipedia.org/wiki/C19H18O8 | {{DISPLAYTITLE:C19H18O8}}
The molecular formula C19H18O8 (molar mass: 374.34 g/mol, exact mass: 374.1002 u) may refer to:
Casticin
Chrysosplenetin
Molecular formulas | C19H18O8 | [
"Physics",
"Chemistry"
] | 64 | [
"Molecules",
"Set index articles on molecular formulas",
"Isomerism",
"Molecular formulas",
"Matter"
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26,374,518 | https://en.wikipedia.org/wiki/Generic%20flatness | In algebraic geometry and commutative algebra, the theorems of generic flatness and generic freeness state that under certain hypotheses, a sheaf of modules on a scheme is flat or free. They are due to Alexander Grothendieck.
Generic flatness states that if Y is an integral locally noetherian scheme, is a finite type morphism of schemes, and F is a coherent OX-module, then there is a non-empty open subset U of Y such that the restriction of F to u−1(U) is flat over U.
Because Y is integral, U is a dense open subset of Y. This can be applied to deduce a variant of generic flatness which is true when the base is not integral. Suppose that S is a noetherian scheme, is a finite type morphism, and F is a coherent OX module. Then there exists a partition of S into locally closed subsets S1, ..., Sn with the following property: Give each Si its reduced scheme structure, denote by Xi the fiber product , and denote by Fi the restriction ; then each Fi is flat.
Generic freeness
Generic flatness is a consequence of the generic freeness lemma. Generic freeness states that if A is a noetherian integral domain, B is a finite type A-algebra, and M is a finite type B-module, then there exists a non-zero element f of A such that Mf is a free Af-module. Generic freeness can be extended to the graded situation: If B is graded by the natural numbers, A acts in degree zero, and M is a graded B-module, then f may be chosen such that each graded component of Mf is free.
Generic freeness is proved using Grothendieck's technique of dévissage. Another version of generic freeness can be proved using Noether's normalization lemma.
References
Bibliography
Algebraic geometry
Commutative algebra
Theorems in abstract algebra | Generic flatness | [
"Mathematics"
] | 413 | [
"Theorems in algebra",
"Fields of abstract algebra",
"Algebraic geometry",
"Commutative algebra",
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23,483,761 | https://en.wikipedia.org/wiki/Molecular%20Cell | Molecular Cell is a peer-reviewed scientific journal that covers research on cell biology at the molecular level, with an emphasis on new mechanistic insights. It was established in 1997 and is published two times per month. The journal is published by Cell Press and is a companion to Cell.
Abstracting and indexing
The journal is abstracted and indexed, for example, in:
According to the Journal Citation Reports, the journal had an impact factor of 14.5 in 2023.
References
External links
Academic journals established in 1997
Molecular and cellular biology journals
Biweekly journals
English-language journals
Cell Press academic journals | Molecular Cell | [
"Chemistry"
] | 123 | [
"Molecular and cellular biology journals",
"Molecular biology"
] |
23,483,982 | https://en.wikipedia.org/wiki/Closed%20range%20theorem | In the mathematical theory of Banach spaces, the closed range theorem gives necessary and sufficient conditions for a closed densely defined operator to have closed range.
The theorem was proved by Stefan Banach in his 1932 Théorie des opérations linéaires.
Statement
Let and be Banach spaces, a closed linear operator whose domain is dense in and the transpose of . The theorem asserts that the following conditions are equivalent:
the range of is closed in
the range of is closed in the dual of
Where and are the null space of and , respectively.
Note that there is always an inclusion , because if and , then . Likewise, there is an inclusion . So the non-trivial part of the above theorem is the opposite inclusion in the final two bullets.
Corollaries
Several corollaries are immediate from the theorem. For instance, a densely defined closed operator as above has if and only if the transpose has a continuous inverse. Similarly, if and only if has a continuous inverse.
Sketch of proof
Since the graph of T is closed, the proof reduces to the case when is a bounded operator between Banach spaces. Now, factors as . Dually, is
Now, if is closed, then it is Banach and so by the open mapping theorem, is a topological isomorphism. It follows that is an isomorphism and then . (More work is needed for the other implications.)
References
.
Banach spaces
Theorems in functional analysis | Closed range theorem | [
"Mathematics"
] | 289 | [
"Theorems in mathematical analysis",
"Theorems in functional analysis"
] |
23,484,427 | https://en.wikipedia.org/wiki/Pregnenedione | A pregnenedione (singular pregnanediol) is an unsaturated diketone derivative of a pregnane. Examples are budesonide and progesterone (pregn-4-ene-3,20-dione).
References
Progestogens | Pregnenedione | [
"Chemistry"
] | 63 | [
"Pharmacology",
"Pharmacology stubs",
"Medicinal chemistry stubs"
] |
31,029,322 | https://en.wikipedia.org/wiki/Pseudo-Boolean%20function | In mathematics and optimization, a pseudo-Boolean function is a function of the form
where is a Boolean domain and is a nonnegative integer called the arity of the function. A Boolean function is then a special case, where the values are also restricted to 0 or 1.
Representations
Any pseudo-Boolean function can be written uniquely as a multi-linear polynomial:
The degree of the pseudo-Boolean function is simply the degree of the polynomial in this representation.
In many settings (e.g., in Fourier analysis of pseudo-Boolean functions), a pseudo-Boolean function is viewed as a function that maps to . Again in this case we can uniquely write as a multi-linear polynomial:
where are Fourier coefficients of and .
Optimization
Minimizing (or, equivalently, maximizing) a pseudo-Boolean function is NP-hard. This can easily be seen by formulating, for example, the maximum cut problem as maximizing a pseudo-Boolean function.
Submodularity
The submodular set functions can be viewed as a special class of pseudo-Boolean functions, which is equivalent to the condition
This is an important class of pseudo-boolean functions, because they can be minimized in polynomial time. Note that minimization of a submodular function is a polynomially solvable problem independent on the presentation form, for e.g. pesudo-Boolean polynomials, opposite to maximization of a submodular function which is NP-hard, Alexander Schrijver (2000).
Roof Duality
If f is a quadratic polynomial, a concept called roof duality can be used to obtain a lower bound for its minimum value. Roof duality may also provide a partial assignment of the variables, indicating some of the values of a minimizer to the polynomial. Several different methods of obtaining lower bounds were developed only to later be shown to be equivalent to what is now called roof duality.
Quadratizations
If the degree of f is greater than 2, one can always employ reductions to obtain an equivalent quadratic problem with additional variables. One possible reduction is
There are other possibilities, for example,
Different reductions lead to different results. Take for example the following cubic polynomial:
Using the first reduction followed by roof duality, we obtain a lower bound of -3 and no indication on how to assign the three variables. Using the second reduction, we obtain the (tight) lower bound of -2 and the optimal assignment of every variable (which is ).
Polynomial Compression Algorithms
Consider a pseudo-Boolean function as a mapping from to . Then Assume that each coefficient is integral. Then for an integer the problem P of deciding whether is more or equal to is NP-complete. It is proved in that in polynomial time we can either solve P or reduce the number of variables to Let be the degree of the above multi-linear polynomial for . Then proved that in polynomial time we can either solve P or reduce the number of variables to .
See also
Boolean function
Quadratic pseudo-Boolean optimization
Notes
References
Mathematical optimization | Pseudo-Boolean function | [
"Mathematics"
] | 631 | [
"Mathematical optimization",
"Mathematical analysis"
] |
31,032,136 | https://en.wikipedia.org/wiki/Engineering%20for%20Change | Engineering for Change (E4C) is an online platform and international community of engineers, scientists, non-governmental organizations, local community advocates and other innovators working to solve problems in sustainable global development. Their mission is to 'prepare, educate, and activate the international technical workforce to improve the quality of life of people and the planet.'
The organization's founding partners are the American Society of Mechanical Engineers, the Institute of Electrical and Electronics Engineers, and Engineers Without Borders USA. It is now under the umbrella of ASME's Engineering for Global Development program. Collaborators include Siemens Stiftung, The Level Market, Autodesk Foundation, Global Alliance for Clean Cookstoves, CAWST, WFEO, ITU, Institute of Food Technologists, and United Nations Major Group for Children and Youth. E4C facilitates the development of affordable, locally appropriate and sustainable solutions to the most pressing humanitarian challenges and shares them freely online as a form of open source appropriate technology.
Members of the E4C community use the platform's online tools to share knowledge, research global development issues, products and services, and deepen their professional development. The organization provides services through seven channels:
The Solutions Library, a database of products that meet basic needs
Monthly webinars and academic seminars
Fellowship Program
News and analysis
Research in collaboration with external partners
Online courses in global development engineering and design
Jobs and volunteer opportunities board
Information about products and services fall into eight categories on the organization's Web site, and they can include big infrastructural projects such as community water purification and bridge building, or smaller, personal technologies such as bicycle-powered electricity generators and cellphone applications for healthcare.
History
In 2009, the American Society of Mechanical Engineers created a website to pull together the disparate sources of information on appropriate technology and solutions in global development. The site aggregated information, hosted a library of often little-known technologies, and offered tools to enable collaboration among development teams worldwide. Throughout 2010, the site operated in alpha and then beta with a mostly closed group of users. A public site, at engineeringforchange.info, mirrored some of the content on the test site, but without all of its functionality. IEEE and EWB-USA signed on as partners in time for the public launch on January 4, 2011.
At present, the organization has more than 70,000 members and a social media following of more than one million.
Fellowship Program
E4C trains and manages ~50 Fellows each year, working with businesses, universities, non-profits and government organizations around the world to improve their programs with technical expertise. Autodesk Foundation is an ongoing supporter of the fellowship, for example. Fellows research, optimize, and design projects that advance the United Nations' Sustainable Development Goals.
Impact Projects
The organization trains and matches technical experts located in dozens of countries to support research, design and other short-term projects in collaboration with other organizations. Impact Projects fall into one of three categories: Design, Research or Workflows. E4C provides management, mentorship, training, and oversight through their Fellowship program.
Solutions Library
The Solutions Library is a database of more than 1,000 products and services that meet basic needs in underserved communities. Entries are investigated and posted by the organization's staff and Fellows worldwide. The information on each entry is standardized to allow users to compare similar products to facilitate research and decision-making.
Education
Education is an important part of Engineering for Change. The Web site provides educational webinars and research seminars, as well as news, analysis, expert opinion and materials on how to design and implement solutions.
See also
Appropriate technology
Open source appropriate technology
References
External links
Official website
"Connecting technical professionals with global design and engineering nonprofits and startups" - Autodesk Foundation
"Innovate for Impact Design Challenge" - Siemens
"Global Coalition Launches Amazon CoLab to Develop and Test Technologies to Protect the Amazon Rainforest" - Conservation X Labs
"A Design Tool Whose Time Has Come" - NextBillion
Innovative Banana Leaf Sanitary Pads Hit a Design Snag - Good.is
American engineering organizations
American Society of Mechanical Engineers
Appropriate technology organizations
Charities based in New York City
Development charities based in the United States
Engineering organizations
Engineers Without Borders
Institute of Electrical and Electronics Engineers
Mechanical engineering organizations
Organizations established in 2009 | Engineering for Change | [
"Engineering"
] | 877 | [
"Engineers Without Borders",
"Institute of Electrical and Electronics Engineers",
"Mechanical engineering organizations",
"nan",
"Mechanical engineering",
"American Society of Mechanical Engineers",
"Electrical engineering organizations"
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41,851,479 | https://en.wikipedia.org/wiki/Pwo%20DNA%20polymerase | Pwo polymerase is a thermostable DNA polymerase used for the polymerase chain reaction. The abbreviation stands for Pyrococcus woesei, a thermophilic archaeon, from which this polymerase was isolated. This polymerase breaks when reaching erroneous uracil in DNA from the chain extension and, through this readahead function, fewer defective DNA clones are synthesized. It is used much less than the usual Taq or Pfu polymerases. This DNA polymerase, similar to other DNA polymerases from Archaebacteria is sensitive to Uracil residues in DNA and is strongly inhibited by dUTP or uracil residues in DNA. Other polymerases in this class are Pfu, Vent, Deep Vent and Pfx. The inhibition of this class of thermostable DNA polymerases limit their use in some applications of PCR, i.e. use of dUTP for prevention of carryover contamination as well as application involving dU containing primers such as ligase free cloning methods or site directed mutagenesis using UNG.
References
DNA replication
EC 2.7.7
Polymerase chain reaction | Pwo DNA polymerase | [
"Chemistry",
"Biology"
] | 246 | [
"Biochemistry methods",
"Genetics techniques",
"Polymerase chain reaction",
"DNA replication",
"Molecular genetics"
] |
41,851,667 | https://en.wikipedia.org/wiki/Vent%20DNA%20polymerase | Vent polymerase is a archean thermostable DNA polymerase used for the polymerase chain reaction. It was isolated from the thermophile Thermococcus litoralis.
References
DNA replication
EC 2.7.7
Polymerase chain reaction | Vent DNA polymerase | [
"Chemistry",
"Biology"
] | 55 | [
"Biochemistry methods",
"Genetics techniques",
"Polymerase chain reaction",
"DNA replication",
"Molecular genetics"
] |
41,854,117 | https://en.wikipedia.org/wiki/Monacolin%20J | Monacolin J is a statin made by red yeast rice. Monacolin J is a precursor to simvastatin and has potential neuroprotective activities.
It can be produced by total mycosynthesis.
References
Statins | Monacolin J | [
"Chemistry"
] | 50 | [
"Organic compounds",
"Organic compound stubs",
"Organic chemistry stubs"
] |
41,858,625 | https://en.wikipedia.org/wiki/Cyclic%20sieving | In combinatorial mathematics, cyclic sieving is a phenomenon in which an integer polynomial evaluated at certain roots of unity counts the rotational symmetries of a finite set.
Given a family of cyclic sieving phenomena, the polynomials give a q-analog for the enumeration of the finite sets, and often arise from an underlying algebraic structure associated to the family of finite sets, such as a representation.
The first study of cyclic sieving was published by Reiner, Stanton and White in 2004. The phenomenon generalizes the "q = −1 phenomenon" of John Stembridge, which considers evaluations of the polynomial only at the first and second roots of unity (that is, q = 1 and q = −1).
Definition
For every positive integer , let denote the primitive th root of unity .
Let be a finite set with an action of the cyclic group , and let be an integer polynomial. The triple exhibits the cyclic sieving phenomenon (or CSP) if for every positive integer dividing , the number of elements in fixed by the action of the subgroup of is equal to . If acts as rotation by , this counts elements in with -fold rotational symmetry.
Equivalently, suppose is a bijection on such that , where is the identity map. Then induces an action of on , where a given generator of acts by . Then exhibits the cyclic sieving phenomenon if the number of elements in fixed by is equal to for every integer .
Example
Let be the set of pairs of elements from . Define a bijection which increases each element in the pair by one (and sends back to ). This induces an action of on , which has an orbit
of size two and an orbit
of size four. If , then counts all elements in , counts fixed points of , counts fixed points of , and counts fixed points of . Hence, the triple exhibits the cyclic sieving phenomenon.
More generally, set and define the q-binomial coefficient by
which is an integer polynomial evaluating to the usual binomial coefficient at . For any positive integer dividing ,
If is the set of size- subsets of with acting by increasing each element in the subset by one (and sending back to ), and if is the q-binomial coefficient above, then exhibits the cyclic sieving phenomenon for every .
In representation theory
The cyclic sieving phenomenon can be naturally stated in the language of representation theory. The group action of on is linearly extended to obtain a representation, and the decomposition of this representation into irreducibles determines the required coefficients of the polynomial .
Let be the vector space over the complex numbers with a basis indexed by a finite set . If the cyclic group acts on , then linearly extending each action turns into a representation of .
For a generator of , the linear extension of its action on gives a permutation matrix , and the trace of counts the elements of fixed by . In particular, the triple exhibits the cyclic sieving phenomenon if and only if for every , where is the character of .
This gives a method for determining . For every integer , let be the one-dimensional representation of in which acts as scalar multiplication by . For an integer polynomial , the triple exhibits the cyclic sieving phenomenon if and only if
Further examples
Let be a finite set of words of the form where each letter is an integer and is closed under permutation (that is, if is in , then so is any anagram of ). The major index of a word is the sum of all indices such that , and is denoted .
If acts on by rotating the letters of each word, and
then exhibits the cyclic sieving phenomenon.
Let be a partition of size with rectangular shape, and let be the set of standard Young tableaux with shape . Jeu de taquin promotion gives an action of on . Let be the following q-analog of the hook length formula:
Then exhibits the cyclic sieving phenomenon. If is the character for the irreducible representation of the symmetric group associated to , then for every , where is the long cycle .
If is the set of semistandard Young tableaux of shape with entries in , then promotion gives an action of the cyclic group on . Define and
where is the Schur polynomial. Then exhibits the cyclic sieving phenomenon.
If is the set of non-crossing (1,2)-configurations of , then acts on these by rotation. Let be the following q-analog of the th Catalan number:
Then exhibits the cyclic sieving phenomenon.
Let be the set of semi-standard Young tableaux of shape with maximal entry , where entries along each row and column are strictly increasing. If acts on by -promotion and
then exhibits the cyclic sieving phenomenon.
Let be the set of permutations of cycle type with exactly exceedances. Conjugation gives an action of on , and if
then exhibits the cyclic sieving phenomenon.
Notes and references
Combinatorics
Generating functions | Cyclic sieving | [
"Mathematics"
] | 996 | [
"Sequences and series",
"Discrete mathematics",
"Mathematical structures",
"Combinatorics",
"Generating functions"
] |
24,970,625 | https://en.wikipedia.org/wiki/RNA-based%20evolution | RNA-based evolution is a theory that posits that RNA is not merely an intermediate between Watson and Crick model of the DNA molecule and proteins, but rather a far more dynamic and independent role-player in determining phenotype. By regulating the transcription in DNA sequences, the stability of RNA, and the capability of messenger RNA to be translated, RNA processing events allow for a diverse array of proteins to be synthesized from a single gene. Since RNA processing is heritable, it is subject to natural selection suggested by Darwin and contributes to the evolution and diversity of most eukaryotic organisms.
Role of RNA in conventional evolution
In accordance with the central dogma of molecular biology, RNA passes information between the DNA of a genome and the proteins expressed within an organism. Therefore, from an evolutionary standpoint, a mutation within the DNA bases results in an alteration of the RNA transcripts, which in turn leads to a direct difference in phenotype.
RNA is also believed to have been the genetic material of the first life on Earth. The role of RNA in the origin of life is best supported by the ease of forming RNA from basic chemical building blocks (such as amino acids, sugars, and hydroxyl acids) that were likely present 4 billion years ago. Molecules of RNA have also been shown to effectively self-replicate, catalyze basic reactions, and store heritable information. As life progressed and evolved over time only DNA, which is much more chemically stable than RNA, could support large genomes and eventually took over the role as the major carrier of genetic information.
Single-Stranded RNA can fold into complex structures
Single-stranded RNA molecules can single handedly fold into complex structures. The molecules fold into secondary and tertiary structures by intramolecular base pairing. There is a fine dynamic of disorder and order that facilitate an efficient structure formation. RNA strands form complementary base pairs. These complementary strands of RNA base pair with another strand, which results in a three-dimensional shape from the paired strands folding in on itself. The formation of the secondary structure results from base pairing by hydrogen bonds between the strands, while tertiary structure results from folding of the RNA. The three-dimensional structure consists of grooves and helices. The formation of these complex structure gives reason to suspect that early life could have formed by RNA.
Variability of RNA processing
Research within the past decade has shown that strands of RNA are not merely transcribed from regions of DNA and translated into proteins. Rather RNA has retained some of its former independence from DNA and is subject to a network of processing events that alter the protein expression from that bounded by just the genomic DNA. Processing of RNA influences protein expression by managing the transcription of DNA sequences, the stability of RNA, and the translation of messenger RNA.
Alternative splicing
Splicing is the process by which non-coding regions of RNA are removed. The number and combination of splicing events varies greatly based on differences in transcript sequence and environmental factors. Variation in phenotype caused by alternative splicing is best seen in the sex determination of D. melanogaster. The Tra gene, determinant of sex, in male flies becomes truncated as splicing events fail to remove a stop codon that controls the length of the RNA molecule. In others the stop signal is retained within the final RNA molecule and a functional Tra protein is produced resulting in the female phenotype. Thus, alternative RNA splicing events allow differential phenotypes, regardless of the identity of the coding DNA sequence.
RNA stability
Phenotype may also be determined by the number of RNA molecules, as more RNA transcripts lead to a greater expression of protein. Short tails of repetitive nucleic acids are often added to the ends of RNA molecules in order to prevent degradation, effectively increasing the number of RNA strands able to be translated into protein. During mammalian liver regeneration RNA molecules of growth factors increase in number due to the addition of signaling tails. With more transcripts present the growth factors are produced at a higher rate, aiding the rebuilding process of the organ.
RNA silencing
Silencing of RNA occurs when double stranded RNA molecules are processed by a series of enzymatic reactions, resulting in RNA fragments that degrade complementary RNA sequences. By degrading transcripts, a lower amount of protein products are translated and the phenotype is altered by yet another RNA processing event.
RNA and Protein
In Earth's early developmental history RNA was the primary substance of life. RNA served as a blueprint for genetic material and was the catalyst to multiply said blueprint. Currently RNA acts by forming proteins. protein enzymes carry out catalytic reactions. RNAs are critical in gene expression and that gene expression depends on mRNA, rRNA, and tRNA. There is a relationship between protein and RNAs. This relationship could suggest that there is a mutual transfer of energy or information. In vitro RNA selection experiments have produced RNA that bind tightly to amino acids. It has been shown that the amino acids recognized by the RNA nucleotide sequences had a disproportionately high frequency of codons for said amino acids. There is a possibility that the direct association of amino acids containing specific RNA sequences yielded a limited genetic code.
Evolutionary mechanism
Most RNA processing events work in concert with one another and produce networks of regulating processes that allow a greater variety of proteins to be expressed than those strictly directed by the genome. These RNA processing events can also be passed on from generation to generation via reverse transcription into the genome. Over time, RNA networks that produce the fittest phenotypes will be more likely to be maintained in a population, contributing to evolution. Studies have shown that RNA processing events have especially been critical with the fast phenotypic evolution of vertebrates—large jumps in phenotype explained by changes in RNA processing events. Human genome searches have also revealed RNA processing events that have provided significant “sequence space for more variability”. On the whole, RNA processing expands the possible phenotypes of a given genotype and contributes to the evolution and diversity of life.
RNA virus evolution
RNA virus evolution appears to be facilitated by a high mutation rate caused by the lack of a proofreading mechanism during viral genome replication. In addition to mutation, RNA virus evolution is also facilitated by genetic recombination. Genetic recombination can occur when at least two RNA viral genomes are present in the same host cell and has been studies in numerous RNA viruses. RNA recombination appears to be a major driving force in viral evolution among Picornaviridae ((+)ssRNA) (e.g. poliovirus). In the Retroviridae ((+)ssRNA)(e.g. HIV), damage in the RNA genome appears to be avoided during reverse transcription by strand switching, a form of genetic recombination. Recombination also occurs in the Coronaviridae ((+)ssRNA) (e.g. SARS). Recombination in RNA viruses appears to be an adaptation for coping with genome damage. Recombination can occur infrequently between animal viruses of the same species but of divergent lineages. The resulting recombinant viruses may sometimes cause an outbreak of infection in humans.
See also
RNA world
References
Evolutionary biology
RNA | RNA-based evolution | [
"Biology"
] | 1,489 | [
"Evolutionary biology"
] |
24,972,374 | https://en.wikipedia.org/wiki/Viscotherm | Viscotherms is a general name used to describe equipment for control of viscosity and temperature of a fluid, in particular of fuel oil in fuel viscosity control systems. The term originated from a brand name Viscotherm registered by VAF Instruments in 1971 and produced until 2009.
The first viscotherm used a measuring capillary and a small pump being fitted throughout the piping of the fluid flow system. The pump produced a constant pressure to force the measured liquid through the capillary. The viscosity was then determined by taking pressure readings a certain points off the capillary. The value of viscosity was then transferred to a PID controller to regulate the temperature of the fuel. This device allowed to measure the real time changes of viscosity and was heavily used in the maritime industry to optimize the combustion of the fuel oil in diesel engines.
After a fire at VAF’s premises in 2009 the Viscotherm was declared obsolete. It was replaced by the Viscosense, a new generation viscometer based on rotational vibrations of a pendulum in a liquid. However, VAF made sure that an easy to install retrofit for the Viscotherm, combining an identical shaped house and the Viscosense sensor, is available.
See also
Viscometer
Fuel oil
Fuel viscosity control
References
H. D. McGeorge. Marine Auxiliary Machinery, Butterworth Heinemann, 1999. ,
External links
VAF Instruments Viscosity Measurement
Viscosity | Viscotherm | [
"Physics"
] | 313 | [
"Physical phenomena",
"Physical quantities",
"Wikipedia categories named after physical quantities",
"Viscosity",
"Physical properties"
] |
24,973,931 | https://en.wikipedia.org/wiki/Extreme%20response%20spectrum | The Extreme Response Spectrum (ERS) (or Maximum Response Spectrum (MRS)) is defined as a curve giving the value of the highest peak of the response of a linear Single Degree of Freedom System (SDOF system) to vibration, according to its natural frequency, for a given damping ratio. The response is described here by the relative movement of the mass of this system in relation to its support. The x-axis refers to the natural frequency and the y-axis to the highest peak multiplied by the square of the quantity (2 π x natural frequency), by analogy with the relative displacement shock response spectrum.
The severity of a vibration can be evaluated by calculating the stresses on a mathematical or finite element model of the structure and, for example, comparison with the ultimate stress of the material. This is the method used to dimension the structure. Generally, however, the problem is instead to evaluate the relative severity of several vibrations (vibrations measured in the real environment, measured vibrations with respect to standards, establishment of a specification etc.). This comparison would be difficult to carry out if one used a fine model of the structure and besides, this is not always available, in particular at the stage of the development of the specification of dimensioning.
A solution consists of applying the vibration under consideration to a “standard” mechanical system, which thus does not claim to be a model of the real structure, composed of a support and N linear one-degree-of-freedom resonators, each one comprising a mass, a spring and a damping device.
A vibration A is considered as more severe than a vibration B if it produces a highest relative displacement (i.e. a highest stress) on this SDOF system than the vibration B.
An ERS is generated from a vibration signal using the following process:
1. Choose a damping ratio for the ERS to be based on;
2. Assume a hypothetical Single Degree of Freedom System, with a given natural frequency (Hz);
3. Calculate (by time base simulation or from a Power Spectral Density (PSD) of the vibratory signal) the highest instantaneous relative displacement experienced by the mass element of this SDOFs at any time during exposure to the vibration in question. Plot this value multiplied by the square of (2 π x natural frequency) against the natural frequency of the hypothetical system;
4. Repeat steps 2 and 3 for other values of the natural frequency.
The resulting plot is called an Extreme response spectrum.
Note
Vibrations can damage a mechanical system as a result of several processes, among which are:
the exceeding of characteristic instantaneous stress limits (yield stress, ultimate stress etc.)
the damage by fatigue following the application of a large number of cycles.
ERS is used according to the first criterion. The second is considered with the fatigue damage spectrum (FDS).
References
Lalanne, C., Mechanical Vibration and Shock Analysis. Volume 5: Specification Development, Second Edition, ISTE - Wiley, 2009.
AFNOR, NF X50-144: Demonstration of resistance to environmental conditions, 2019
Mechanical vibrations | Extreme response spectrum | [
"Physics",
"Engineering"
] | 627 | [
"Structural engineering",
"Mechanics",
"Mechanical vibrations"
] |
24,974,609 | https://en.wikipedia.org/wiki/Florida%20Platform | The Florida Platform is a flat geological feature with the emergent portion forming the Florida peninsula.
Structure
The platform forms an escarpment between the Gulf of Mexico and the Atlantic Ocean. The platform's western edge, or Florida Escarpment, is normally defined where water depths at drop dramatically and in a short distance to . The Florida peninsula is located on the eastern side of the platform, where in places it lies only from the platform's edge. On the gulf side the platform ends over to the west of the modern shoreline, where a massive cliff rises over from the depth of the gulf floor. The western reaches of the platform just off Tampa were explored by the submersible DSV Alvin. Examination has placed the depth of carbonate rocks at greater than .
Geology
The platform's basement rocks include Precambrian–Cambrian igneous rocks, Ordovician–Devonian sedimentary rocks, and Triassic–Jurassic volcanic rock. Florida's igneous and sedimentary foundation separated from what is now the African Plate when the supercontinent Pangaea rifted apart in the middle Triassic and possibly early-middle Jurassic. It then secured to the North American craton. One of the early geologic structures is the Peninsular Arch which formed during the Jurassic.
The oldest sediments that are exposed are Eocene carbonates found in the Avon Park Formation. Most of the state of Florida is covered by Pliocene, Pleistocene, and Holocene siliciclastic-bearing sediments deposited during sea-level fluctuations and filling in of the Gulf Trough beginning in the late Tertiary and Quaternary.
Notes
External links
Geologic Map of the State of Florida
General Facts about the Gulf of Mexico
Geology of Florida
Physical oceanography
Terranes
Landforms of the Gulf of Mexico | Florida Platform | [
"Physics"
] | 353 | [
"Applied and interdisciplinary physics",
"Physical oceanography"
] |
24,976,623 | https://en.wikipedia.org/wiki/Hydrogen%20pinch | Hydrogen pinch analysis (HPA) is a hydrogen management method that originates from the concept of heat pinch analysis. HPA is a systematic technique for reducing hydrogen consumption and hydrogen generation through integration of hydrogen-using activities or processes in the petrochemical industry, petroleum refineries hydrogen distribution networks and hydrogen purification.
Principle
A mass analysis is done by representing the purity and flowrate for each stream from the hydrogen consumers (sinks), such as hydrotreaters, hydrocrackers, isomerization units and lubricant plants and the hydrogen producers (sources), such as hydrogen plants and naphtha reformers, streams from hydrogen purifiers, membrane reactors, pressure swing adsorption and continuous distillation and off-gas streams from low- or high-pressure separators. The source-demand diagram shows bottlenecks, surplus or shortages. The hydrogen pinch is the purity at which the hydrogen network has neither hydrogen surplus nor deficit.
After the analysis REFOPT from the Centre for Process Integration at The University of Manchester is used as a tool for process integration with which the process is optimized. The methodology was also developed into commercial software by companies such as Linnhoff March and AspenTech. The Aspen product incorporated the work of Nick Hallale (formerly a lecturer at University of Manchester) and was the first method to consider multiple components, rather than a pseudo-binary mixture of hydrogen and methane.
History
The first assessment based on cost and value composite curves of hydrogen resources of a hydrogen network was proposed by Tower et al. (1996). Alves developed the hydrogen pinch analysis approach based on the concept of heat pinch analysis in 1999. Nick Hallale and Fang Liu extended this original work, adding pressure constraints and mathematical programming for optimisation. This was followed by developments at AspenTech, producing commercial software for industrial application.
See also
Water Pinch
Timeline of hydrogen technologies
Review of process integration by Nick Hallale (University of Manchester) covering hydrogen pinch
References
Citations
Sources
Nick Hallale, Ian Moore, Dennis Vauk, "Hydrogen optimization at minimal investment", Petroleum Technology Quarterly (PTQ), Spring (2003)
External links
Hydrogen Pinch made easy
Mechanical engineering
Chemical engineering
Hydrogen technologies
Analysis | Hydrogen pinch | [
"Physics",
"Chemistry",
"Engineering"
] | 450 | [
"Chemical engineering",
"nan",
"Applied and interdisciplinary physics",
"Mechanical engineering"
] |
24,979,028 | https://en.wikipedia.org/wiki/Concrete%20degradation | Concrete degradation may have many different causes. Concrete is mostly damaged by the corrosion of reinforcement bars due to the carbonatation of hardened cement paste or chloride attack under wet conditions. Chemical damage is caused by the formation of expansive products produced by chemical reactions (from carbonatation, chlorides, sulfates and distillate water), by aggressive chemical species present in groundwater and seawater (chlorides, sulfates, magnesium ions), or by microorganisms (bacteria, fungi...) Other damaging processes can also involve calcium leaching by water infiltration, physical phenomena initiating cracks formation and propagation, fire or radiant heat, aggregate expansion, sea water effects, leaching, and erosion by fast-flowing water.
The most destructive agent of concrete structures and components is probably water. Indeed, water often directly participates to chemical reactions as a reagent and is always necessary as a solvent, or a reacting medium, making transport of solutes and reactions possible. Without water, many harmful reactions cannot progress, or are so slow that their harmful consequences become negligible during the planned service life of the construction. Dry concrete has a much longer lifetime than water saturated concrete in contact with circulating water. So, when possible, concrete must first be protected from water infiltrations.
Corrosion of reinforcement bars
The expansion of the corrosion products (iron oxides) of carbon steel reinforcement structures may induce internal mechanical stress (tensile stress) that cause the formation of cracks and disrupt the concrete structure. If rebars have been improperly installed or have inadequate concrete cover at surfaces exposed to the elements, oxide jacking and spalling can occur during the structure's lifetime: flat fragments of concrete are detached from the concrete mass as a result of the rebar's corrosion.
Concrete, like most consolidated hard rocks, is a material very resistant to compression but which cannot withstand tension, especially internal tensions. Its tensile strength is about 10 times lower than its compressive strength. In itself carbonated concrete is a very solid material because its compressive strength increases due to its porosity decrease by the precipitation of calcium carbonate (calcite, ). In the absence of steel reinforcement bars and without the formation of expansive reaction products inducing tensile stress inside the concrete matrix, pure concrete is most often a long-lasting material. An illustration of the concrete intrinsic durability is the dome of the Pantheon building in Rome made with Roman concrete more than 2000 years ago.
When atmospheric carbon dioxide (), or carbonate ions (, dissolved in water) diffuse into concrete from its external surface, they react with calcium hydroxide (portlandite, ) and the pH of the concrete pore water progressively decreases from 13.5 – 12.5 to 8.5 (pH of water in equilibrium with calcite). Below a pH value of about 9.5 – 10, the solubility of iron oxides present at the surface of carbon steel increases and they start to dissolve. As a consequence, they no longer protect the underlying metallic iron against oxidation by atmospheric oxygen and the reinforcement bars are no longer passivated against corrosion. It is the considerable forces internally created by the expansion of the iron corrosion products (about 6 – 7 times less dense than metallic iron, so 6 – 7 times more voluminous) that cause the cracks in the concrete matrix and destroy reinforced concrete. In the absence of iron (and without some harmful chemical degradation reactions also producing expansive products) concrete would probably be one of the most durable materials. However, steel reinforcement bars are necessary to take over the tensile efforts to which concrete is submitted in most engineering structures and stainless steel would be too costly a metal to replace carbon steel. Zinc-galvanization or epoxy-coating can improve the corrosion resistance of rebar, but have also other disadvantages such as their lower surface adhesion to concrete (risk of slip), the possible formation of cathodic and anodic zones conducive to galvanic corrosion if the protective coating is locally punctured or damaged, and their higher costs.
Formation of expansive phases in concrete
As hard rocks, concrete can withstand high compressive stress but not tensile stress. As a consequence, concrete is easily damaged when expansion phases are formed in its mass.
The most ubiquitous, and best known, expansive phases are probably the iron oxides produced by the oxidation of the carbon steel reinforcement bars embedded into concrete. Corrosion products are formed around rebar located in carbonated concrete (and thus no longer passivated against corrosion), or directly exposed to the atmospheric oxygen once cracks have started to form. The damages produced by rebar corrosion are clearly visible with bare eyes and their diagnostic is easy.
Other deleterious expansive chemical reactions more difficult to characterize and to identify can occur in concrete. They may be first distinguished according to the location where they occur in concrete: inside the aggregates or in the hardened cement paste.
Expansion inside the aggregates
Various types of aggregates can undergo different chemical reactions and swell inside concrete, leading to damaging expansive phenomena.
Alkali–silica reaction
The most common are those containing reactive amorphous silica, that can react in the presence of water with the cement alkalis (K2O and Na2O). Among the more reactive siliceous mineral components of some aggregates are opal, chalcedony, flint and strained quartz. Silica (in fact silicic acid when hydrated) is easily dissolved by sodium hydroxide (NaOH) to form sodium silicate (), a strong desiccant with a high affinity for water. This reaction is at the core of the alkali–silica reaction (ASR):
Following this reaction, a hygroscopic and expansive viscous silicagel phase forms inside the affected aggregates which swell and crack from inside. In its turn, the volumetric expansion of the swollen aggregates damages the concrete matrix and extensive cracks propagate causing structural damages in the concrete structure. On the surface of concrete pavements, the ASR can also cause pop-outs, i.e. the expulsion of small cones (up to in diameter), corresponding to aggregate particle size.
A quite similar reaction (alkali-silicate reaction) can occur when clay minerals are present in some impure aggregates, and it may also lead to destructive expansion.
Alkali–carbonate reaction
With some aggregates containing dolomite, a dedolomitization reaction, also known as alkali-carbonate reaction (ACR), can occur where the magnesium carbonate () reacts with the hydroxyl ions () and yields magnesium hydroxide (brucite, ) and a carbonate ion (). The resulting expansion caused by the swelling of brucite can cause destruction of the material:
Often the alkali–silicate reaction and the dedolomitization reaction are masked by a much more severe alkali–silica reaction dominating the deleterious effects. Because the alkali-carbonate reaction (ACR) is often thwarted by a coexisting ASR reaction, it explains why ACR is no longer considered to be a major detrimental reaction.
Pyrite oxidation
Far less common are degradation and pop-outs caused by the presence of pyrite (), a disulfide (S-S) very sensitive to oxidation by atmospheric oxygen, that generates expansion by forming less dense insoluble iron oxides (), iron oxy-hydroxides (FeO(OH), or ·n) and mildly soluble gypsum (·2).
When complete (i.e., when all ions are also oxidized into less soluble ions), pyrite oxidation can be globally written as follows:
The sulfuric acid released by pyrite oxidation then reacts with portlandite ()) present in the hardened cement paste to give gypsum:
When concrete is carbonated by atmospheric carbon dioxide (), or if limestone aggregates are used in concrete, reacts with calcite () and water to also form gypsum while releasing back to the atmosphere:
The dihydrated gypsum is relatively soluble in water at room temperature and thus mobile. It can easily be leached by infiltration water and can form efflorescences on the concrete surface while the insoluble ·n remain in place around the grains of oxidized pyrite they taint in red-ocre.
Expansive chemical reactions inside the hardened cement paste
The sulfate anions reacting with different phases of the hardened cement paste (HCP) to form more voluminous reaction products can cause 3 types of expansive reactions called sulfate attack inside HCP:
The delayed ettringite formation (DEF) also known as internal sulfate attack (ISA);
The external sulfate attack (ESA), and;
The thaumasite form of sulfate attack (TSA).
These three types of sulfate attack reactions are described into more details in specific sections latter in the text. When the hardened cement paste (HCP) is affected, the detrimental consequences for the structural stability of concrete structures are generally more severe than when aggregates are affected: DEF, ESA and TSA are much more damaging for concrete than ASR and ACR reactions.
A common points to all these various chemical expansive reaction is that they all require water as a reactant and as a reaction medium. The presence of water is always an aggravating factor. Concrete structures immersed in water as dams and bridge piles are therefore particularly sensitive. These reactions are also characterized by slow reaction kinetics, depending on environmental conditions such as temperature and relative humidity. They develop at a slow rate and may take several years before damages become apparent. Often a decade is needed to observe their harmful consequences. Protecting concrete structures from water contact may help to slow down the progression of the damages.
Chemical damages
Carbonation
Carbon dioxide () from air and bicarbonate () or carbonate () anions dissolved in water react with the calcium hydroxide (, portlandite) produced by Portland cement hydration in concrete to form calcium carbonate () while releasing a water molecule in the following reaction:
Exception made of the water molecule, the carbonation reaction is essentially the reverse of the process of calcination of limestone taking place in a cement kiln:
Carbonation of concrete is a slow and continuous process of atmospheric diffusing from the outer surface of concrete exposed to air into its mass and chemically reacting with the mineral phases of the hydrated cement paste. Carbonation slows down with increasing diffusion depth.
Carbonation has two antagonist effects for (1) the concrete strength, and (2) its durability:
The precipitation of calcite filling the microscopic voids in the concrete pore space decreases the concrete matrix porosity: so, it increases the mechanical strength of concrete;
At the same time carbonation consumes portlandite and therefore decreases the concrete alkalinity reserve buffer. Hyper-alkaline conditions (i.e., basic chemical conditions) characterized by a high pH (typically 12.5 – 13.5) are needed to passivate the steel surface of the reinforcement bars (rebar) and to protect them from corrosion. Below a pH of 10, the solubility of the iron oxides forming a protective thin coating at the surface of carbon steel increases. The thin protective oxide layer starts to dissolve, and corrosion is then promoted. As the volumetric mass of iron oxides can be as high as 6 – 7 times that of metallic iron (Fe), a detrimental consequence is the expansion of the corrosion products around the rebar. This causes the development of a tensile stress in the concrete matrix around the rebar. When the tensile strength of concrete is exceeded in the concrete cover above the rebar, concrete starts to spall. Cracks appear in the concrete cover protecting the rebar against corrosion and constitute preferential pathways for direct ingress towards the rebar. This accelerates the carbonation reaction and in turn the corrosion process speeds up.
This explain why the carbonation reaction of reinforced concrete is an undesirable process in concrete chemistry. Concrete carbonation can be visually revealed by applying a phenolphthalein solution over the fresh surface of a concrete samples (concrete core, prism, freshly fractured bar). Phenolphthalein is a pH indicator, whose color turns from colorless at pH < 8.5 to pink-fuchsia at pH > 9.5. A violet color indicates still alkaline areas and thus non-carbonated concrete. Carbonated zones favorable for steel corrosion and concrete degradation are colorless.
The presence of water in carbonated concrete is necessary to lower the pH of concrete pore water around rebar and to depassivate the carbon steel surface at low pH. Water is central to corrosion processes. Without water, the steel corrosion is very limited and rebar present in dry carbonated concrete structures, or components, not affected by water infiltration do not suffer from significant corrosion.
Chloride attack
The main effect of chloride ions on reinforced concrete is to cause pitting corrosion of the steel reinforcement bars (rebar). It is a surreptitious and dangerous form of localized corrosion because the rebar sections can be decreased to the point that the steel reinforcement are no longer capable to withstand to the tensile efforts they are supposed to resist by design. When the rebar sections are too small or the rebar are locally broken, the reinforcements are lost, and concrete is no longer reinforced concrete.
Chlorides, particularly calcium chloride, have been used to shorten the setting time of concrete.
However, calcium chloride and (to a lesser extent) sodium chloride have been shown to leach calcium hydroxide and cause chemical changes in Portland cement, leading to loss of strength, as well as attacking the steel reinforcement present in most concrete. The ten-storey Queen Elizabeth hospital in Kota Kinabalu contained a high percentage of chloride causing early failure.
Alkali–silica reaction (ASR)
The alkali–silica reaction (ASR) is a deleterious chemical reaction between the alkali ( and ), dissolved in concrete pore water as NaOH and KOH, with reactive amorphous (non-crystalline) siliceous aggregates in the presence of moisture. The simplest way to write the reaction in a stylized manner is the following (other representations also exist):
(young N-S-H gel)
This reaction produces a gel-like substance of sodium silicate ( • n ), also noted • n , or N-S-H (sodium silicate hydrate). This hygroscopic gel swells inside the affected reactive aggregates which expand and crack. In its turn, it causes concrete expansion. If concrete is heavily reinforced, it can first cause some prestressing effect before cracking and damaging the structure. ASR affects the aggregates and is recognizable by cracked aggregates. It does not directly affect the hardened cement paste (HCP).
Delayed ettringite formation (DEF, or ISA)
When the temperature of concrete exceeds 65 °C for too long a time at an early age, the crystallization of ettringite (AFt) does not occur because of its higher solubility at elevated temperature and the then less soluble mono-sulfate (AFm) is formed. After dissipation of the cement hydration heat, temperature goes back to ambient and the temperature curves of the solubilities of AFt and AFm phases cross over. The mono-sulfate (AFm) now more soluble at low temperature slowly dissolves to recrystallize as the less soluble ettringite (AFt). AFt crystal structure hosts more water molecules than AFm. So, AFt has a higher molar volume than AFm because of its 32 H2O molecules. During months, or years, after young concrete cooling, AFt crystallizes very slowly as small acicular needles and can exert a considerable crystallization pressure on the surrounding hardened cement paste (HCP). This leads to the expansion of concrete, to its cracking, and it can ultimately lead to the ruin of the affected structure. The characteristic feature of delayed ettringite formation (DEF) is a random honeycomb cracking pattern similar to this of the alkali-silica reaction (ASR). In fact, this typical crack pattern is common to all expansive internal reactions and also to restrained shrinkage where a rigid substrate or a dense rebar network prevents the movements of a superficial concrete layer. DEF is also known as internal sulfate attack (ISA). External sulfate attack (ESA) also involves ettringite (AFt) formation and deleterious expansion with the same harmful symptoms but requires an external source of sulfate anions in the surrounding terrains or environment. To avoid DEF or ISA reactions, the best way is to use a low C3A (tri-calcium aluminate) cement precluding the formation of ettringite (AFt). Sulfate resisting (SR) cements have also a low content in Al2O3. DEF, or ISA, only affects the hardened cement paste (HCP) and leaves intact the aggregates.
DEF is exacerbated at high pH in cement with a too high content in alkalis and therefore in hydroxides. This is caused by the transformation of ettringite (AFt) into aluminoferrite monosulfate (AFm) under the action of the hydroxyl anions (OH–) as schematized as follows:
AFt + OH– → AFm
The complete reaction can be derived from the molecular formulas of the reagents and products involved in the reaction. This reaction favors the dissolution of AFt and the formation of AFm. When combined, it is an aggravating factor of the harmful effect of too high temperatures. To minimize DEF, the use of low-alkali cements is also recommended. The detrimental crystallization of ettringite (AFt) preferentially occurs when concrete is exposed to water infiltrations and that the pH decreases due to the leaching of the (OH–) ions: the reaction is reversed as when temperature decreases.
External sulfate attacks (ESA)
Sulfates in solution in contact with concrete can cause chemical changes to the cement, which can cause significant microstructural effects leading to the weakening of the cement binder (chemical sulfate attack). Sulfate solutions can also cause damage to porous cementitious materials through crystallization and recrystallization (salt attack). Sulfates and sulfites are ubiquitous in the natural environment and are present from many sources, including gypsum (calcium sulfate) often present as an additive in 'blended' cements which include fly ash and other sources of sulfate. With the notable exception of barium sulfate, most sulfates are slightly to highly soluble in water. These include acid rain where sulfur dioxide in the airshed is dissolved in rainfall to produce sulfurous acid. In lightning storms, the dioxide is oxidized to trioxide making the residual sulfuric acid in rainfall even more highly acidic. Concrete sewage infrastructure is most commonly attacked by sulfuric acid and sulfate anions arising from the oxidation of sulfide present in the sewage. Sulfides are formed when sulfate-reducing bacteria present in sewer mains reduce the ubiquitous sulfate ions present in water drains into hydrogen sulfide gas (). is volatile and released from water in the sewage atmosphere. It dissolves in a thin film of water condensed onto the wall of the sewer ducts where it is also accompanied by hydrogeno-sulfide () and sulfide () ions. When and anions are further exposed to atmospheric oxygen or to oxygenated stormwater, they are readily oxidized and produce sulfuric acid (in fact acidic hydrogen ions accompanied by sulfate and bisulfate ions) according to the respective oxidation reactions:
or,
The corrosion often present in the crown (top) of concrete sewers is directly attributable to this process – known as crown rot corrosion.
Thaumasite form of sulfate attack (TSA)
Thaumasite is a calcium silicate mineral, containing Si atoms in unusual octahedral configuration, with chemical formula Ca3Si(OH)6(CO3)(SO4)·12H2O, also sometimes more simply written as CaSiO3·CaCO3·CaSO4·15H2O.
Thaumasite is formed under special conditions in the presence of sulfate ions in concrete containing, or exposed to, a source of carbonate anions such as limestone aggregates, or finely milled limestone filler (). Bicarbonate anions () dissolved in groundwater may also contribute to the reaction. The detrimental reaction proceeds at the expense of calcium silicate hydrates (C-S-H, with dashes denoting here their non-stoichiometry) present in the hardened cement paste (HCP). The thaumasite form of sulfate attack (TSA) is a particular type of very destructive sulfate attack. C-S-H are the "glue" in the hardened cement paste filling the interstices between the concrete aggregates. As the TSA reaction consumes the silicates of the "cement glue", it can lead to a harmful decohesion and a softening of concrete. Expansion and cracking are more rarely observed. Unlike the common sulfate attack, in which the calcium hydroxide (portlandite) and calcium aluminate hydrates react with sulfates to respectively form gypsum and ettringite (an expansive phase), in the case of TSA the C-S-H ensuring the cohesion of HCP and aggregates are destroyed. As a consequence, even concrete containing low- sulfate-resisting Portland cement may be affected.
TSA is sometimes easily recognizable on the field when examining the altered concrete. TSA-affected concrete becomes powdery and can be dug with a scoop, or even scrapped with the fingers. Concrete decohesion is very characteristic of TSA.
TSA was first identified during the years 1990 in England in the United Kingdom in the foundation piles of bridges of the motorway M5 located in the Kimmeridgian marls. These marls are a mixture of clay and limestone sedimented under anoxic conditions and are rich in pyrite (, a disulfide). Once these marls were excavated, pyrite was exposed to atmospheric oxygen or oxygen-rich infiltration water and rapidly oxidized. Pyrite oxidation produces sulfuric acid. In its turn, reacts with portlandite (present in the hardened cement paste, HCP) and calcite ( (present in limestone aggregates or in carbonated HCP). The strong acidification of the medium caused by pyrite oxidation releases bicarbonate ions () or carbon dioxide () along with calcium () and sulfate ions ().
Full pyrite oxidation can be schematized as:
The sulfuric acid released by pyrite oxidation then reacts with portlandite ()) present in the hardened cement paste to give gypsum:
When concrete also contains limestone aggregates or a filler addition, reacts with calcite () and water to also form gypsum while releasing :
Gypsum is relatively soluble in water , so there is plenty of calcium and sulfates ions available for TSA.
Simultaneously, carbonic acid () dissolves calcite to form soluble calcium bicarbonate:
So, when all the chemical ingredients necessary to react with C-S-H from the hardened cement paste in concrete are present together the TSA reaction can occur. When grounds rich in pyrite, such as many clays or marls, are excavated for civil engineering works, the strong acidification produced by pyrite oxidation is the powerful driving force triggering TSA because it frees up and mobilizes all the ions needed to attack C-S-H and to form thaumasite ().
TSA is favored by a low temperature, although it can be encountered at higher temperature in warm areas. The reason is to be found in the retrograde solubility of most of the ingredients needed for the TSA reaction. Indeed, the solubility of dissolved carbon dioxide (), portlandite (), calcite (), and gypsum (·2), increases when the temperature is lowered. This is because the dissolution reactions of these mineral species are exothermic and release heat. A lower temperature facilitates the heat release and therefore favors the exothermic reaction. Only the solubility of silica (from C-S-H) increases with temperature because silica dissolution is an endothermic process which requires heat to proceed.
Calcium leaching
When water flows through cracks present in concrete, water may dissolve various minerals present in the hardened cement paste or in the aggregates, if the solution is unsaturated with respect to them. Dissolved ions, such as calcium (Ca2+), are leached out and transported in solution some distance. If the physico-chemical conditions prevailing in the seeping water evolve with distance along the water path and water becomes supersaturated with respect to certain minerals, they can further precipitate, making calthemite deposits (predominately calcium carbonate) inside the cracks, or at the concrete outer surface. This process can cause the self-healing of fractures in particular conditions.
Fagerlund (2000) determined that, “About 15% of the lime has to be dissolved before strength is affected. This corresponds to about 10% of the cement weight, or almost all of the initially formed Ca(OH)2.” Therefore, a large amount of "calcium hydroxide" (Ca(OH)2) must be leached from the concrete before structural integrity is affected. The other issue however is that leaching away Ca(OH)2 may allow the corrosion of reinforcing steel to affect structural integrity.
Decalcification
Within set concrete there remains some free "calcium hydroxide" (Ca(OH)2), which can further dissociate to form Ca2+ and hydroxide (OH−) ions". Any water which finds a seepage path through micro cracks and air voids present in concrete, will readily carry the (Ca(OH)2) and Ca2+ (depending on solution pH and chemical reaction at the time) to the underside of the structure where leachate solution contacts the atmosphere. Carbon dioxide (CO2) from the atmosphere readily diffuses into the leachate and causes a chemical reaction, which precipitates (deposits) calcium carbonate (CaCO3) on the outside of the concrete structure. Consisting primarily of CaCO3 this secondary deposit derived from concrete is known as "calthemite" and can mimic the shapes and forms of cave "speleothems", such as stalactites, stalagmites, flowstone etc. Other trace elements such as iron from rusting reinforcing steel bars may be transported and deposited by the leachate at the same time as the CaCO3. This may colour the calthemites orange or red.
The chemistry involving the leaching of calcium hydroxide from concrete can facilitate the growth of calthemites up to ≈200 times faster than cave speleothems due to the different chemical reactions involved. The sight of calthemite is a visual sign that calcium is being leached from the concrete structure and the concrete is gradually degrading.
In very old concrete where the calcium hydroxide has been leached from the leachate seepage path, the chemistry may revert to that similar to "speleothem" chemistry in limestone cave. This is where carbon dioxide enriched rain or seepage water forms a weak carbonic acid, which leaches calcium carbonate (CaCO3) from within the concrete structure and carries it to the underside of the structure. When it contacts the atmosphere, carbon dioxide degasses and calcium carbonate is precipitated to create calthemite deposits, which mimic the shapes and forms of speleothems. This degassing chemistry is not common in concrete structures as the leachate can often find new paths through the concrete to access free calcium hydroxide and this reverts the chemistry to that previously mentioned where CO2 is the reactant.
Sea water attack
Concrete exposed to seawater is susceptible to its corrosive effects. The effects are more pronounced above the tidal zone than where the concrete is permanently submerged. In the submerged zone, magnesium and hydrogen carbonate ions precipitate a layer of brucite (magnesium hydroxide: Mg(OH)2), about 30 micrometers thick, on which a slower deposition of calcium carbonate as aragonite occurs. These mineral layers somewhat protect the concrete from other processes, which includes attack by magnesium, chloride and sulfate ions and carbonatation. Above the water surface, mechanical damage may occur by erosion by waves themselves or sand and gravel they carry, and by crystallization of salts from water soaking into the concrete pores and then drying up. Pozzolanic cements and cements using more than of blast furnace slags as cementitious material are more resistant to seawater than pure Portland cement. Seawater attack presents aspects of both chloride and sulfate attacks.
Effects of bacterial activity
Bacteria themselves do not have noticeable effect on concrete. However, sulfate-reducing bacteria (SRB) in untreated sewage water tend to produce hydrogen sulfide (H2S), which is then oxidized in sulfuric acid (H2SO4) by atmospheric oxygen (abiotic reaction) and by aerobic bacteria present in biofilm (biotic reaction) on the concrete surface above the water level. The sulfuric acid dissolves the carbonates in the hardened cement paste (HCP), and also calcium hydroxide (portlandite: Ca(OH)2) and calcium silicate hydrate (CaO·SiO2·nH2O), and causes strength loss, as well as producing sulfates which are harmful to concrete.
H2SO4 + Ca(OH)2 → CaSO4 + 2 H2O
H2SO4 + CaO·SiO2·n H2O → CaSO4 + H2SiO3 + n H2O
In each case the soft expansive and water-soluble corrosion product of gypsum (CaSO4) is formed. Gypsum is easily washed away in wastewater causing a loss of concrete aggregate and exposing fresh material to acid attack.
Concrete floors lying on ground that contains pyrite (iron(II) disulfide) are also at risk.
As a preventive measure sewage may be pretreated to increase pH or oxidize or precipitate the sulfides in order to minimize the activity of sulfide-reducing bacteria.
As bacteria often prefer to adhere to the surfaces of solids than to remain into suspension in water (planktonic bacteria), the biofilms formed by sessile (i.e., fixed) bacteria are often the place where they are the most active. Biofilms made of multiple layers (like an onion) of dead and living bacteria protect the living ones from the harsh conditions often prevailing in water outside biofilm. Biofilms developing on the already exposed surface of metallic elements encased in concrete can also contribute to accelerate their corrosion (differential aeration and formation of anodic zones at the surface of the metal). Sulfides produced by the SRB bacteria can also induce stress corrosion cracking in steel and other metals.
Physical damages
Construction defects
Damages can occur during the casting and de-shuttering processes. For instance, the corners of beams can be damaged during the removal of shuttering because they are less effectively compacted by means of vibration (improved by using form-vibrators). Other physical damages can be caused by the use of steel shuttering without base plates. The steel shuttering pinches the top surface of a concrete slab due to the weight of the next slab being constructed.
Concrete slabs, block walls and pipelines are susceptible to cracking during ground settlement, seismic tremors or other sources of vibration, and also from expansion and contraction during adverse temperature changes.
Various types of concrete shrinkage
Chemical shrinkage (self-desiccation)
The cement hydration process consumes water molecules. The sum of the volumes of the hydration products present in the hardened cement paste is smaller than the sum of the volumes of the reacting mineral phases present in the cement clinker. Therefore, the volume of the fresh and very young concrete undergoes a contraction due to the hydration reaction: it is what is called "chemical shrinkage" or "self-desiccation". It is not a problem as long as the very fresh concrete is still in a liquid, or a sufficiently plastic, state and can easily accommodate volume changes (contraction).
Plastic shrinkage
Later in the setting phase, when the fresh concrete becomes more viscous and starts to harden, water loss due to unwanted evaporation can cause "plastic shrinkage". This occur when concrete is placed under hot conditions, e.g. in the summer and not sufficiently protected against evaporation. Cracks often develop above reinforcement bars because the contraction of concrete is locally restrained at this level and the still setting and weakly resistant concrete cannot freely shrink.
Cracks due to a poor curing (loss of water at early age)
The curing of concrete when it continues to harden after its initial setting and progressively develops its mechanical strength is a critical phase to avoid unwanted cracks in concrete. Depending on the temperature (summer or winter conditions) and thus on the cement hydration kinetics controlling the setting and hardening rate of concrete, curing time can require a few days only (summer) or up to two weeks (winter). It is then capital to avoid losses of water by evaporation because water is still necessary for continuing the slow cement hydration. Water loss at this stage aggravates concrete shrinkage and can cause unacceptable cracks to develop in concrete. Cracks form in case of a too short, or too poor, curing when young concrete has not yet developed a sufficient early strength to withstand tensile stress caused by undesirable and premature drying. Cracks development occurs when early-age concrete is insufficiently protected against desiccation and too much water evaporates with heat because of unfavorable meteorological conditions: e.g, high temperature, direct solar insolation, dry air, low relative humidity, and high wind speed during summer, or in hot conditions. Curing is intended to maintain moist conditions at the surface of concrete. It can be done by letting the formworks in place for a longer time, or by applying a hydrophobic thin film of an oily product (curing compound) at the concrete surface (e.g., for large slabs or rafts) to minimize water evaporation.
Drying shrinkage
After sufficient setting and hardening of concrete (after 28 days), the progressive loss of capillary water is also responsible for the "drying shrinkage". It is a continuous and long-term process occurring later during the concrete life when under dry conditions the larger pores of concrete are no longer completely saturated by water.
Thermal cracks
When concrete is subject to an excessive temperature increase during its setting and hardening as in massive concrete structures from where cement hydration heat cannot easily escape (semi-adiabatic conditions), the temperature gradients and the differential volume changes can also cause the formation of thermal cracks and fissures. To minimize them a slowly-setting cement (CEM III, with blast furnace slags) is preferred to a quickly setting cement (CEM I: Portland cement). Pouring concrete under colder conditions (e.g., during the night, or in the winter), or using cold water and ice mixed with cooled aggregates to prepare concrete, may also contribute to minimize thermal cracks.
Restrained shrinkage
When a concrete structure is heavily reinforced, the very dense rebar network can block the contraction movement of the protecting concrete cover located above the external layer of reinforcement bars due to the natural drying shrinkage process. As a consequence, a network of fissures with the characteristic honeycomb pattern also typical for the cracks resulting from the expansive chemical reactions (ASR, DEF, ESA) forms.
The formation of fissures in the concrete cover above the reinforcement bars represents a preferential pathway for the ingress of water and aggressive agents such as (lowering of pH around the rebar) and chloride anions (pitting corrosion) into concrete. The physical formation of cracks therefore favors the chemical degradation of concrete and aggravates steel corrosion. Physical and chemical degradation processes are intimately coupled, and the presence of water infiltrations also accelerates the formation of expansive products of harmful swelling chemical reactions (iron corrosion products, ASR, DEF, ISA, ESA).
Different approaches and methods have been developed to attempt to quantitatively estimate the influence of cracks in concrete structures on carbonation and chloride penetration. Their aim is to avoid underestimating the penetration depth of these harmful chemical agents and to calculate a sufficient thickness for the concrete cover to protect the rebar against corrosion during the whole service life of the concrete structure.
Freeze-thaw cycles
In winter conditions, or in cold climates, when the temperature falls below , the crystallization of ice in the pores of concrete is also a physical mechanism (change of state) responsible for the volumetric expansion of a substance exerting a high tensile strength inside the concrete matrix. When the tensile strength of concrete is exceeded, cracks appear. Adding an air entrainment agent during the mixing of fresh concrete induces the formation of tiny air bubbles in the fresh concrete slurry. This creates numerous small air-filled micro-cavities in the hardened concrete serving as empty volume reserve to accommodate the volumetric expansion of ice and delays the moment tensile stress will develop. Air entrainment makes concrete more workable during placement, and increases its durability when hardened, particularly in climates subject to freeze-thaw cycles.
Mechanical damages
Overload, shocks and vibrations (bridges, roads submitted to intense truck traffic...) can induce mechanical stress and deformations in concrete structures and be responsible for the mechanical degradation of concrete. Beside the long-term drying shrinkage of concrete, pre-stressed and post-tensioned civil engineering structures (bridges, primary containment domes of nuclear power plants can also undergo slow concrete creep and deformation.
Thermal damages
Due to its low thermal conductivity, a layer of concrete is frequently used for fireproofing of steel structures. However, concrete itself may be damaged by fire, with one notable example being the 1996 Channel Tunnel fire where fire damage extended along several hundred meters of the tunnel's length. For this reason, common fire testing standards, such as ASTM E119, do not permit fire testing of cementitious products unless the relative humidity inside the cementitious product is at or below 75%. Otherwise, concrete can be subject to significant spalling.
Up to about 300 °C, the concrete undergoes normal thermal expansion. Above that temperature, shrinkage occurs due to water loss; however, the aggregate continues expanding, which causes internal stresses. Up to about 500 °C, the major structural changes are carbonatation and coarsening of pores. At 573 °C, quartz undergoes rapid expansion due to phase transition, and at 900 °C calcite starts shrinking due to decomposition. At 450-550 °C the cement hydrate decomposes, yielding calcium oxide. Calcium carbonate decomposes at about 600 °C. Rehydration of the calcium oxide on cooling of the structure causes expansion, which can cause damage to material which withstood fire without falling apart. Concrete in buildings that experienced a fire and were left standing for several years shows extensive degree of carbonatation from carbon dioxide which is reabsorbed.
Concrete exposed to up to 100 °C is normally considered as healthy. The parts of a concrete structure that is exposed to temperatures above approximately 300 °C (dependent of water/cement ratio) will most likely get a pink color. Over approximately 600 °C the concrete will turn light grey, and over approximately 1000 °C it turns yellow-brown.
One rule of thumb is to consider all pink colored concrete as damaged that should be removed.
Fire will expose the concrete to gases and liquids that can be harmful to the concrete, among other salts and acids that occur when gases produced by a fire come into contact with water.
If concrete is exposed to very high temperatures very rapidly, explosive spalling of the concrete can result. In a very hot, very quick fire the water inside the concrete will boil before it evaporates. The steam inside the concrete exerts expansive pressure and can initiate and forcibly expel a spall.
Radiation damages
Exposure of concrete structures to neutrons and gamma radiation in nuclear power plants, and high-flux material testing reactors, can induce radiation damage to their concrete structural components. Paramagnetic defects and optical centers are easily formed, but very high fluxes are necessary to displace a sufficiently high number of atoms in the crystal lattice of the minerals present in concrete before significant mechanical damage is observed.
However, neutron irradiation with a very high neutron fluence (number of neutrons per unit of cross-section area: neutron/cm2) is known to render amorphous a fraction of the quartz present in some concrete aggregates. This amorphization process is also called metamictization. Metamict quartz with its disordered lattice structure is prone to alkali–silica reaction and can thus be responsible of harmful chemical expansion in the concrete of nuclear containment structures.
Repairs and strengthening
It may be necessary to repair a concrete structure following damage (e.g. due to age, chemical attack, fire, impact, movement or reinforcement corrosion). Strengthening may be necessary if the structure is weakened (e.g. due to design or construction errors, excessive loading, or because of a change of use).
Repair techniques
The first step should always be an investigation to determine the cause of the deterioration. The general principles of repair include arresting and preventing further degradation; treating exposed steel reinforcement; and filling fissures or holes caused by cracking or left after the loss of spalled or damaged concrete.
Various techniques are available for the repair, protection and rehabilitation of concrete structures, and specifications for repair principals have been defined systematically. The selection of the appropriate approach will depend on the cause of the initial damage (e.g. impact, excessive loading, movement, corrosion of the reinforcement, chemical attack, or fire) and whether the repair is to be fully load bearing or simply cosmetic.
Concrete stitching employs metal staples or stitches to restore structural integrity to cracked concrete surfaces. This method applies torque across the crack, effectively transferring load and tension to stabilize and strengthen the affected area. Recognized for its simplicity and minimal disruption, concrete stitching is widely utilized in essential infrastructures such as bridges and buildings, significantly prolonging the lifespan of concrete structures by preventing crack propagation.
Repair principles which do not improve the strength or performance of concrete beyond its original (undamaged) condition include replacement and restoration of concrete after spalling and delamination; strengthening to restore structural load-bearing capacity; and increasing resistance to physical or mechanical attack.
Repair principles for arresting and preventing further degradation include control of anodic areas; cathodic protection, cathodic control; increasing resistivity; preserving or restoring passivity; increasing resistance to chemical attack; protection against ingress of adverse agents; and moisture control.
Techniques for filling holes left by the removal of spalled or damaged concrete include mortar repairs; flowing concrete repairs and sprayed concrete repairs. The filling of cracks, fissures or voids in concrete for structural purposes (restoration of strength and load-bearing capability), or non-structural reasons (flexible repairs where further movement is expected, or alternately to resist water and gas permeation) typically involves the injection of low viscosity resins or grouts based on epoxy, PU or acrylic resins, or micronised cement slurries.
One novel proposal for the repair of cracks is to use bacteria. BacillaFilla is a genetically engineered bacterium designed to repair damaged concrete, filling in the cracks, and making them whole again.
Strengthening techniques
Various techniques are available for strengthening concrete structures, to increase the load-carrying capacity or else to improve the in-service performance. These include increasing the concrete cross-section and adding material such as steel plate or fiber composites to enhance the tensile capacity or increase the confinement of the concrete for improved compression capacity.
See also
Calthemite
Concrete fracture analysis
Corrosion of rebar
Electrical resistivity measurement of concrete
Interfacial Transition Zone (ITZ)
Pitting corrosion of rebar
Reinforced concrete structures durability
References
Further reading
Cement
Concrete degradation | Concrete degradation | [
"Materials_science"
] | 9,204 | [
"Materials degradation",
"Concrete degradation"
] |
24,980,042 | https://en.wikipedia.org/wiki/Amanda%20Barnard | Amanda Susan Barnard (born 31 December 1971) is an Australian theoretical physicist working in predicting the real world behavior of nanoparticles using analytical models and supercomputer simulations and applied machine learning. Barnard is a pioneer in the thermodynamic cartography of nanomaterials, creating nanoscale phase diagrams relevant to different environmental conditions, and relating these to structure/property maps. Her current research involves developing and applying statistical methods and machine/deep learning in nanoscience and nanotechnology, and materials and molecular informatics. In 2014 she became the first person in the southern hemisphere, and the first woman, to win the Feynman Prize in Nanotechnology, which she won for her work on diamond nanoparticles.
Barnard is currently based in Australia as Professor of Computational Science in the Research School of Computer Science at the Australian National University.
Biography
In 2001, she graduated with a first-class honours science degree from the Royal Melbourne Institute of Technology (RMIT), majoring in applied physics. Barnard received a PhD in 2003 from RMIT for her computer modelling work predicting and explaining various forms of nanocarbon at different sizes. Following her PhD, Barnard served as a Distinguished Postdoctoral Fellow in the Center for Nanoscale Materials at Argonne National Laboratory (USA). She also held a senior research position as Violette & Samuel Glasstone Fellow at the University of Oxford (UK) with an Extraordinary Research Fellowship at The Queen's College. Professor Barnard then moved to CSIRO as an ARC Queen Elizabeth II Fellow, an Office of the Chief Executive Science Leader, and finally as a Chief Research Scientist spanning 2009 to 2020.
Qualifications
2003 Doctor of Philosophy (Physics), RMIT University
2001 Bachelor of Science, First Class Honours (Applied Physics), RMIT University
Career highlights, awards, fellowships and grants
2022 Member of the Order of Australia (AM) in the 2022 Australia Day Honours for "significant service to computational science, to medical research, and to education".
2019 AMMA Medal, Association of Molecular Modellers of Australasia
2017 Woman of Achievement, Black & White Foundation, Australia
2014 Feynman Prize in Nanotechnology (Theory)
2014 ACS Nano Lectureship (Asia/Pacific), American Chemical Society, USA
2010 IEEE Distinguished Lecturer Award, IEEE, South Australia
2010 UNSW Eureka Prize for Scientific Research, Australian Museum
2010 Frederick White Prize, Australian Academy of Science
2009 Malcolm McIntosh Prize for Physical Scientist of the Year
2009– Leader of the Virtual Nanoscience Laboratory, CSIRO Materials Science and Engineering
2009— Queen Elizabeth II Fellowship, Australian Research Council
2009 Mercedes-Benz Australian Environmental Research Award, Banksia Environmental Foundation
2009 Young Scientist Prize in Computational Physics, International Union of Pure and Applied Physics
2009 JG Russell Award, Australian Academy of Science
2009 Future Summit Leadership Award, Australian Davos Connection
2008 L'Oréal Australia For Women in Science Fellowship
2008 Alumnus of the Year, RMIT University
2008 Inaugural Future Generation Fellowship, School of Chemistry, University of Melbourne
2005–2008 Extraordinary Junior Research Fellowship, Queen's College, Oxford, UK
2005–2008 Violette & Samuel Glasstone Fellowship, Department of Materials, University of Oxford, UK
2004 Innovation Award (Student Category), RMIT University
2004 University Research Prize, RMIT University
2003–2005 Distinguished Postdoctoral Fellowship, Center for Nanoscale Materials, Argonne National Laboratory, USA
Research highlights
Identified the link between nanomorphology and the environmental stability of nanomaterials, and how it influences reactivity and potential "nano-hazards"
Developed a new technique for investigating the shape of nanomaterials as a function of size, temperature or chemical potential, able to include experimentally realistic structures and chemical environments
First researcher to report investigations into the effect of shape on size-dependent phase transitions in nanomaterials
Discovered the first example of anisotropic (facet-dependent) surface electrostatic potential in a homoelemental nanomaterial, resulting in dipolar or multipolar interactions in a non-polar material
Leader in statistical nanoscience and the use of statistical analysis and machine learning to predict the properties of diverse and complex ensembles of nanoscale materials.
References
External links
CSIRO staff profile
2008 L'Oréal For Women in Science Australian Fellowships award citation
21st-century Australian physicists
Living people
1971 births
Computational physicists
Australian women physicists
Members of the Order of Australia
RMIT University alumni | Amanda Barnard | [
"Physics"
] | 894 | [
"Computational physicists",
"Computational physics"
] |
24,980,412 | https://en.wikipedia.org/wiki/Thalidomide%20scandal | In the late 1950s and early 1960s, the use of thalidomide in 46 countries was prescribed to women who were pregnant or who subsequently became pregnant, and consequently resulted in the "biggest anthropogenic medical disaster ever," with more than 10,000 children born with a range of severe deformities, such as phocomelia, as well as thousands of miscarriages.
Thalidomide was introduced in 1953 as a tranquilizer, and was later marketed by the German pharmaceutical company Chemie Grünenthal under the trade name Contergan as a medication for anxiety, trouble sleeping, tension, and morning sickness. It was introduced as a sedative and medication for morning sickness without having been tested on pregnant women. While initially deemed to be safe in pregnancy, concerns regarding birth defects were noted in 1961, and the medication was removed from the market in Europe that year.
Development of thalidomide
Thalidomide was first developed as a tranquilizer by Swiss pharmaceutical company Ciba in 1953.
In 1954, Ciba abandoned the product, and it was acquired by German pharmaceutical company Chemie Grünenthal. The company had been established by Hermann Wirtz Sr, a Nazi Party member, after World War II as a subsidiary of the family's Mäurer & Wirtz company. The company's initial aim was to develop antibiotics for which there was an urgent market need. Wirtz appointed chemist Heinrich Mückter, who had escaped prosecution for war crimes for his experiments on prisoners of Nazi concentration camps, to head the development programme because of his experience researching and producing an anti-typhus vaccine for Nazi Germany. He hired Martin Staemmler, a medical doctor and leading proponent of the Nazi eugenics programme, as head of pathology, as well as Heinz Baumkötter, the chief medical officer at the Sachsenhausen concentration camp, and Otto Ambros, a chemist and Nazi war criminal. Ambros was the chairman of Grünenthal's advisory committee during the development of thalidomide and was a board member when Contergan was being sold.
Birth defect crisis
The total number of embryos affected by the use of thalidomide during pregnancy is estimated at more than 10,000, and potentially up to 20,000; of these, approximately 40 percent died at or shortly after the time of birth. Those who survived had limb, eye, urinary tract, and heart defects. Its initial entry into the U.S. market was prevented by Frances Oldham Kelsey at the U.S. Food and Drug Administration (FDA). The birth defects of thalidomide led to the development of greater drug regulation and monitoring in many countries.
The severity and location of the deformities depended on how many days into the pregnancy the mother was before beginning treatment; thalidomide taken on the 20th day of pregnancy caused central brain damage, day 21 would damage the eyes, day 22 the ears and face, day 24 the arms, and leg damage would occur if taken up to day 28. Thalidomide did not damage the fetus if taken after 42 days' gestation.
United Kingdom
In the UK, the drug was licensed in 1958 and withdrawn in 1961. Of the approximately 2,000 babies born with defects, around half died within a few months and 466 survived to at least 2010. In 1968, after a long campaign by The Sunday Times, a compensation settlement for the UK victims was reached with Distillers Company (now part of Diageo), which had distributed the drug in the UK. Distillers Biochemicals paid out approximately £28m in compensation following a legal battle.
The British Thalidomide Children's Trust was set up in 1973 as part of a £20 million legal settlement between Distillers Company and 429 children with thalidomide-related disabilities. In 1997, Diageo (formed by a merger between Grand Metropolitan and Guinness, who had taken over Distillers in 1990) made a long-term financial commitment to support the Thalidomide Trust and its beneficiaries. The UK government gave survivors a grant of £20 million, to be distributed through the Thalidomide Trust, in December 2009.
Spain
In Spain, thalidomide was widely available throughout the 1970s, and perhaps even into the 1980s. There were two reasons for this. First, state controls and safeguarding were poor; it was not until 2008 that the government even admitted the country had ever imported thalidomide. Second, Grünenthal failed to insist that its sister company in Madrid warn Spanish doctors, and permitted its sister company not to warn doctors of the defects. The Spanish advocacy group for victims of thalidomide estimates that in 2015, there were 250–300 living victims of thalidomide in Spain.
Australia and New Zealand
Although the Australian obstetrician William McBride took credit for raising concern about thalidomide, it was a midwife called Sister Pat Sparrow who first suspected the drug was causing birth defects in the babies of patients under McBride's care at Crown Street Women's Hospital in Sydney. German paediatrician Widukind Lenz, who also suspected the link, is credited with conducting the scientific research that proved thalidomide was causing birth defects in 1961. McBride was later awarded a number of honors, including a medal and prize money by L'Institut de la Vie in Paris, but he was eventually struck off the Australian medical register in 1993 for scientific fraud related to work on Debendox. Further animal tests were conducted by George Somers, Chief Pharmacologist of Distillers Company in Britain, which showed fetal abnormalities in rabbits. Similar results were also published showing these effects in rats and other species.
Lynette Rowe, who was born without limbs, led an Australian class action lawsuit against the drug's manufacturer, Grünenthal, which fought to have the case heard in Germany. The Supreme Court of Victoria dismissed Grünenthal's application in 2012, and the case was heard in Australia. On 17 July 2012, Rowe was awarded an out-of-court settlement, believed to be in the millions of dollars and providing precedence for class action victims to receive further compensation. In February 2014, the Supreme Court of Victoria endorsed the settlement of $89 million AUD to 107 victims of the drug in Australia and New Zealand.
Germany
In East Germany, thalidomide was rejected by the Central Committee of Experts for the Drug Traffic in the GDR, and was never approved for use. There are no known thalidomide babies born in East Germany. Meanwhile, in West Germany, it took some time before the increase in dysmelia at the end of the 1950s was connected with thalidomide. In 1958, Karl Beck, a former pediatric doctor in Bayreuth, wrote an article in a local newspaper claiming a relationship between nuclear weapons testing and cases of dysmelia in children. Based on this, FDP whip Erich Mende requested an official statement from the federal government. For statistical reasons, the main data series used to research dysmelia cases started by chance at the same time as the approval date for thalidomide. After the Nazi regime with its Law for the Prevention of Hereditarily Diseased Offspring used mandatory statistical monitoring to commit various crimes, western Germany had been very reluctant to monitor congenital disorders in a similarly strict way. The parliamentary report rejected any relation with radioactivity and the abnormal increase of dysmelia. Also the DFG research project installed after the Mende request was not helpful. The project was led by pathologist Franz Büchner, who ran the project to propagate his teratological theory. Büchner saw lack of healthy nutrition and behavior of the mothers as being more important than genetic reasons. Furthermore, it took a while to appoint a Surgeon General in Germany; the Federal Ministry of Health was not founded until 1962, some months after thalidomide was banned from the market. In West Germany approximately 2,500 babies were born with birth defects from thalidomide.
Canada
Despite its severe side effects, thalidomide was sold in pharmacies in Canada until 1962. The effects of thalidomide increased fears regarding the safety of pharmaceutical drugs. The Society of Toxicology of Canada was formed after the effects of thalidomide were made public, focusing on toxicology as a discipline separate from pharmacology. The need for the testing and approval of the toxins in certain pharmaceutical drugs became more important after the disaster. The Society of Toxicology of Canada is responsible for the Conservation Environment Protection Act, focusing on researching the impact to human health of chemical substances. Thalidomide brought on changes in the way drugs are tested, what type of drugs are used during pregnancy, and increased the awareness of potential side effects of drugs.
According to Canadian news magazine programme W5, most, but not all, victims of thalidomide receive annual benefits as compensation from the Government of Canada. Excluded are those who cannot provide the documentation the government requires.
A group of 120 Canadian survivors formed the Thalidomide Victims Association of Canada, the goal of which is to prevent the approval of drugs that could be harmful to pregnant individuals and babies. The members from the thalidomide victims association were involved in the STEPS programme, which aimed to prevent teratogenicity.
United States
In the U.S., the FDA refused approval to market thalidomide, saying further studies were needed. This reduced the impact of thalidomide in U.S. patients. The refusal was largely due to pharmacologist Frances Oldham Kelsey who withstood pressure from the Richardson-Merrell Pharmaceuticals Co. Although thalidomide was not approved for sale in the United States at the time, over 2.5 million tablets had been distributed to over 1,000 physicians during a clinical testing programme. It is estimated that nearly 20,000 patients, several hundred of whom were pregnant, were given the drug to help alleviate morning sickness or as a sedative, and at least 17 children were consequently born in the United States with thalidomide-associated deformities. While pregnant, children's television host Sherri Finkbine took thalidomide that her husband had purchased over-the-counter in Europe. When she learned that thalidomide was causing fetal deformities she wanted to abort her pregnancy, but the laws of Arizona allowed abortion only if the mother's life was in danger. Finkbine traveled to Sweden to have the abortion. Thalidomide was found to have deformed the fetus.
For denying the application despite the pressure from Richardson-Merrell Pharmaceuticals Co., Kelsey eventually received the President's Award for Distinguished Federal Civilian Service at a 1962 ceremony with President John F. Kennedy. In September 2010, the FDA honored Kelsey with the first Kelsey award, given annually to an FDA staff member. This came 50 years after Kelsey, then a new medical officer at the agency, first reviewed the application from the William S. Merrell Pharmaceuticals Company of Cincinnati.
Cardiologist Helen B. Taussig learned of the damaging effects of the drug thalidomide on newborns and in 1967, testified before Congress on this matter after a trip to Germany where she worked with infants with phocomelia (severe limb deformities). As a result of her efforts, thalidomide was banned in the United States and Europe.
Austria
Ingeborg Eichler, a member of the Austrian pharmaceutical admission conference, enforced restrictions on the sale of thalidomide (tradename Softenon) under the rules of prescription medication and as a result relatively few affected children were born in Austria and Switzerland.
Aftermath of scandal
The numerous reports of malformations in babies brought about the awareness of the side effects of the drug on pregnant women. The birth defects caused by the drug thalidomide can range from moderate malformation to more severe forms. Possible birth defects include phocomelia, dysmelia, amelia, bone hypoplasticity, and other congenital defects affecting the ear, heart, or internal organs. Franks et al. looked at how the drug affected newborn babies, the severity of their deformities, and reviewed the drug in its early years. Webb in 1963 also reviewed the history of the drug and the different forms of birth defects it had caused. "The most common form of birth defects from thalidomide is shortened limbs, with the arms being more frequently affected. This syndrome is the presence of deformities of the long bones of the limbs resulting in shortening and other abnormalities."
Grünenthal criminal trial
In 1968, a large criminal trial began in West Germany, charging several Grünenthal officials with negligent homicide and injury. After Grünenthal settled with the victims in April 1970, the trial ended in December 1970 with no finding of guilt. As part of the settlement, Grünenthal paid 100 million DM into a special foundation; the West German government added 320 million DM. The foundation paid victims a one-time sum of 2,500–25,000 DM (depending on severity of disability) and a monthly stipend of 100–450 DM. The monthly stipends have since been raised substantially and are now paid entirely by the government (as the foundation had run out of money). Grünenthal paid another €50 million into the foundation in 2008.
On 31 August 2012, Grünenthal chief executive Harald F. Stockwho served as the chief executive officer of Grünenthal GmbH from January 2009 to May 28, 2013apologized for the first time for producing the drug and remaining silent about the birth defects. At a ceremony, Stock unveiled a statue of a disabled child to symbolize those harmed by thalidomide and apologized for not trying to reach out to victims for over 50 years. At the time of the apology, there were between 5,000 and 6,000 people still living with Thalidomide-related birth defects. Victim advocates called the apology "insulting" and "too little, too late", and criticized the company for not compensating victims and for their claim that no one could have known the harm the drug caused, arguing that there were plenty of red flags at the time.
Australian National Memorial
On 13 November 2023, the Australian Government announced its intention to make a formal apology to people affected by thalidomide with the unveiling of a national memorial site. Prime Minister Anthony Albanese described the thalidomide tragedy as a “dark chapter” in Australian history, and Health Minister Mark Butler said, “While we cannot change the past or end the physical suffering, I hope these important next steps of recognition and apology will help heal some of the emotional wounds.”
Notable cases
Mercédes Benegbi, born with phocomelia of both arms, drove the successful campaign for compensation from her government for Canadians who were affected by thalidomide.
Mat Fraser, born with phocomelia of both arms, is an English rock musician, actor, writer and performance artist. He produced a 2002 television documentary, Born Freak, which looked at this historical tradition and its relevance to modern disabled performers. This work has become the subject of academic analysis in the field of disability studies.
Niko von Glasow, a thalidomide survivor, produced a documentary called NoBody's Perfect, based on the lives of 12 people affected by the drug, which was released in 2008.
Josée Lake is a Canadian Paralympic gold medallist swimmer, thalidomide survivor, and president of the Thalidomide Victims Association of Canada
Lorraine Mercer MBE of the United Kingdom, born with phocomelia of both arms and legs, is the only thalidomide survivor to carry the Olympic Torch.
Thomas Quasthoff, an internationally acclaimed bass-baritone, who describes himself: "1.34 meters tall, short arms, seven fingers — four right, three left — large, relatively well-formed head, brown eyes, distinctive lips; profession: singer".
Alvin Law, Canadian motivational speaker and former radio broadcaster.
Change in drug regulations
The disaster prompted many countries to introduce tougher rules for the testing and licensing of drugs, such as the Kefauver Harris Amendment (US), Directive 65/65/EEC1 (E.U.), and the Medicines Act 1968 (UK). In the United States, the new regulations strengthened the FDA, among other ways, by requiring applicants to prove efficacy and to disclose all side effects encountered in testing. The FDA subsequently initiated the Drug Efficacy Study Implementation to reclassify drugs already on the market.
References
Further reading
External links
WHO Pharmaceuticals Newsletter No. 2, 2003 – See page 11, Feature Article
CBC Digital Archives – Thalidomide: Bitter Pills, Broken Promises
Remind me again, what is thalidomide and how did it cause so much harm?. The Conversation, 7 December 2015
Congenital amputations
Drug safety
20th-century health disasters
Health disasters in the United Kingdom
Leprosy
Medical controversies
Medical controversies in the United Kingdom
Medical controversies in Germany
Medical controversies in Austria
Medical controversies in Australia
Medical controversies in New Zealand
Medical controversies in the United States
Medical controversies in Canada
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Wikipedia medicine articles ready to translate
History of disability | Thalidomide scandal | [
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24,981,838 | https://en.wikipedia.org/wiki/Idiopathic%20hypersomnia | Idiopathic hypersomnia (IH) is a neurological disorder which is characterized primarily by excessive sleep and excessive daytime sleepiness (EDS). Idiopathic hypersomnia was first described by Bedrich Roth in 1976, and it can be divided into two forms: polysymptomatic and monosymptomatic. The condition typically becomes evident in early adulthood and most patients diagnosed with IH will have had the disorder for many years prior to their diagnosis. , an FDA-approved medication exists for IH called Xywav, which is an oral solution of calcium, magnesium, potassium, and sodium oxybates; in addition to several off-label treatments (primarily FDA-approved narcolepsy medications).
Idiopathic hypersomnia may also be referred to as IH, IHS, or primary hypersomnia, and belongs to a group of sleep disorders known as central hypersomnias, central disorders of hypersomnolence, or hypersomnia of brain origin. Diagnostic and Statistical Manual of Mental Disorders, Fourth Edition (DSM-IV) defines idiopathic hypersomnia as EDS without narcolepsy or the associated features of other sleep disorders. It occurs in the absence of medical problems or sleep disruptions, such as sleep apnea, that can cause secondary hypersomnia.
Signs and symptoms
Individuals with IH share common symptoms including excessive daytime sleepiness, sleep inertia, brain fog, and long sleep periods.
Excessive daytime sleepiness, characterized by persistent sleepiness throughout the day and often a general lack of energy, even during the day after apparently adequate or even prolonged nighttime sleep. People with EDS nap repeatedly throughout the day and have strong urges to sleep while driving, working, eating, or conversing with others.
Sleep inertia (also known as sleep drunkenness), characterized by having extreme difficulty waking up and feeling an uncontrollable desire to go back to sleep.
Clouding of consciousness characterized by inattention, thought process abnormalities, comprehension abnormalities, and language abnormalities. These symptoms may affect performance on perception, memory, learning, executive functions, language, constructive abilities, voluntary motor control, attention, and mental speed. Affected individuals can complain of forgetfulness, confusion, or inability to think clearly.
Excessive sleep (9 hours or more over a full 24-hour period), without feeling refreshed after waking. Daytime naps can be up to several hours and are also unrefreshing.
Some studies have shown increased frequencies of palpitations, digestive problems, difficulty with body temperature regulation, and other symptoms in patients with IH. Anxiety and depression are frequent, and are likely to be a response to chronic illness. A case series in 2010 found that peripheral vascular symptoms, such as cold hands and feet (e.g., Raynaud syndrome), were more common in people with IH than in controls.
Other autonomic dysfunctional symptoms, such as fainting episodes, dizziness upon arising, possibly migrainous headaches, food cravings, and impotence may also be correlated with IH. Researchers have found that people reporting IH symptoms report high levels of autonomic dysfunction on par with other conditions of autonomic failure (i.e., MSA and diabetes). there is no evidence that such symptoms are related to IH.
Causes
Unlike narcolepsy with cataplexy, which has a known cause (autoimmune destruction of hypocretin-producing neurons), the cause of IH is largely unknown. , researchers have identified a few abnormalities associated with IH, which with further study may help to clarify the etiology.
Destruction of noradrenergic neurons has produced hypersomnia in experimental animal studies, and injury to adrenergic neurons has also been shown to lead to hypersomnia. IH has also been associated with a malfunction of the norepinephrine system and decreased cerebrospinal fluid (CSF) histamine levels.
Researchers have recently found an abnormal hypersensitivity to GABA (the major brain chemical responsible for sedation) in a subset of patients with central hypersomnia i.e.: IH, narcolepsy without cataplexy and long sleepers. They have identified a small (500 to 3000 daltons) naturally occurring bioactive substance (most likely a peptide as it is trypsin-sensitive) in the CSF of affected patients. Although this substance requires further identification of its chemical structure, it is currently referred to as a "somnogen" because it has been shown to cause hyper-reactivity of GABAA receptors, which leads to increased sedation or somnolence. In essence, it is as though these patients are chronically sedated with a benzodiazepine (medication which acts through the GABA system) such as midazolam or alprazolam, even though they do not take these medications.
Diagnosis
Idiopathic hypersomnia lacks a clearly defining biological marker (i.e., HLA-DQB1*0602 genotype in narcolepsy). Doctors can more carefully exclude these causes of EDS in order to more correctly diagnose IH. However, "even in the presence of other specific causes of hypersomnia, one should carefully assess the contribution of these etiological factors to the complaint of EDS and when specific treatments of these conditions fail to suppress EDS, the additional diagnosis of IH should be considered."
The severity of EDS can be quantified by subjective scales, such as the Epworth Sleepiness Scale and the Stanford Sleepiness Scale, and also by objective tests, like actigraphy, psychomotor vigilance task, maintenance of wakefulness test (MWT), multiple sleep latency test (MSLT) although as per latest research studies, the effectiveness of MSLT has been called into question. Several groups of researchers have found normal MSLT results in patients who otherwise seem to have IH. Therefore, when IH is suspected, researchers suggest appending a 24-hour continuous polysomnography to the standard overnight/MSLT study in order to record total sleep time.
The recent studies have also found that reports of sleepiness in IH relate more to mental fatigue rather than physiological sleepiness per se and the subjective scales like ESS, IH symptom diary (IHSD) and PGIC better captures the severity of symptoms consistently.
It is also important to note that whereas narcolepsy is strongly associated with the HLA-DQB1*0602 genotype, "HLA typing is of no help in the positive diagnosis of idiopathic hypersomnia." This is "despite some reports that suggest an increase [sic] frequency of HLA Cw2 and DRS in idiopathic hypersomnia subjects."
In patients with IH, polysomnography typically shows short sleep latency, increased mean slow wave sleep, and a high mean sleep efficiency. "Latency to REM sleep and percentages of light sleep and REM sleep were normal, compared with normal ranges." Despite this, one study has found increased sleep fragmentation in patients with IH without long sleep time, suggesting multiple possible presentations.
Per ICSD-3, five criteria must be met for a diagnosis of IH:
Daytime lapses into sleep or an irrepressible need to sleep on a daily basis, for at least 3 months
Absence of insufficient sleep syndrome
Absence of cataplexy
Absence of other causes of hypersomnia
The presence of positive MSLT tests.
The latest ICD-10 defines IH with long sleep time as a neurological disorder that is a rare sleep disorder characterized by prolonged sleep at night and extreme sleepiness during the day. There are no apparent causes. This disorder affects the ability to function. It is of central nervous system origin characterized by prolonged nocturnal sleep and periods of daytime drowsiness. Affected individuals experience difficulty with awakening in the morning and may have associated sleep drunkenness, automatic behaviors, and memory disturbances. This condition differs from narcolepsy in that daytime sleep periods are longer, there is no association with cataplexy, and the multiple sleep latency onset test does not record sleep-onset rapid eye movement sleep.
Management
Since the underlying mechanism is not yet fully understood, treatment efforts have usually focused on symptom management. In August 2021, low-sodium oxybate (Xywav) became the first U.S. FDA-approved treatment for idiopathic hypersomnia.
Wake-promoting medications used in narcolepsy are also commonly used off-label to help manage the excessive daytime sleepiness of IH. Stimulants tend to be less effective for IH than they are for narcolepsy and may be less well tolerated.
Stimulants
The main treatment of excessive daytime sleepiness is done through central nervous system stimulants.
Methylphenidate and dextroamphetamine are most used stimulants to controlled EDS. Increased dopamine release is felt to be the main property explaining wake-promotion from these medications. Insomnia is another common side effect and may require additional treatment.
Solriamfetol is a dopamine and norepinephrine reuptake inhibitor (NDRI) used to treat excessive daytime sleepiness associated with narcolepsy and obstructive sleep apnea.
Pitolisant, a selective histamine 3 (H3) receptor antagonist/inverse agonist, was approved by FDA during August 2019. It works by increasing the synthesis and release of histamine, a wake-promoting neurotransmitter in the brain.
Modafinil and Armodafinil elevate hypothalamic histamine levels, and they are known to bind to the dopamine transporter, thereby inhibiting dopamine reuptake. Modafinil can cause uncomfortable side effects, including nausea, headache, and a dry mouth for some patients, while other patients report no noticeable improvement even on relatively high dosages. They may also "interact with low-dose contraceptives, potentially reducing efficacy, although the scientific data supporting this claim is weak and rests on poorly documented anecdotes."
Sleep-promoting medications
Sodium oxybate (Xyrem) is an orphan drug which was designed specifically for the treatment of narcolepsy. Common side effects include nausea, dizziness, and hallucinations. A 2016 study by Leu-Semenescu et al. found sodium oxybate reduced daytime sleepiness in IH to the same degree as in patients with narcolepsy type 1, and the drug improved severe sleep inertia in 71% of the hypersomnia patients. In July 2020, the FDA approved Xywav™ (calcium, magnesium, potassium, and sodium oxybates), an oral solution for the treatment of cataplexy or excessive daytime sleepiness (EDS) in patients 7 years of age and older with narcolepsy.
Cognitive behavioral therapy
Although behavioral approaches have not been demonstrated to improve clinical markers of IH, cognitive behavioral therapy (CBT) has been found to potentially reduce depressive symptoms and improve self-efficacy in people with hypersomnia.
Prognosis
IH can profoundly affect work, education, and quality of life due to excessive daytime somnolence. Patients will often need to adapt their lifestyle to avoid situations that might be dangerous while sleepy, such as high risk work, or driving. The risks associated with these activities have been found to be higher for patients with hypersomnias than for those with sleep apnea or severe insomnia.
Epidemiology
Typically, the symptoms of IH begin in adolescence or young adulthood, although they can begin at a later age. After onset, hypersomnia often worsens over several years, but it is often stable by the time of diagnosis and appears to be a lifelong condition. Spontaneous remission is only seen in 10–15% of patients.
According to the limited epidemiological data that exists, IH "has more of a female preponderance (1.8/1)." Family cases are frequent, in a range from 25% to 66% without any clear mode of inheritance."
IH has long been considered a rare disease, believed to be 10 times less frequent than narcolepsy. The prevalence of narcolepsy (with cataplexy) is estimated between 1/3,300 and 1/5,000. Although the true prevalence of IH is unknown, it is estimated at 1/10,000 to 1/25,000 for the long sleep form and 1/11,000 to 1/100,000 without long sleep. A more precise estimate "is complicated by a lack of clear biologic markers" and a lack of "unambiguous diagnostic criteria."
Because of the rarity of IH, research into the condition is limited "Patients are rare, researchers and scientists involved in the field are few and research findings are therefore scarce." "In Europe and in North America there is now a public health concern about helping patients and families affected by these rare diseases. Due to the complexity of the disease, they often experience difficulties to be diagnosed and often face social and professional consequences."
Research
GABA-directed medications
There is ongoing research into the efficacy of gamma aminobutyric acid A (GABAA) receptor antagonists for the treatment of IH. Research findings suggest that the GABA neurotransmitter system plays a significant role in the etiology of primary hypersomnias, such as IH and Narcolepsy Type 2. Given the possible role of hyperactive GABAA receptors in IH, medications that could counteract this activity are being studied to test their potential to improve sleepiness. These currently include clarithromycin and flumazenil.
Flumazenil, a GABAA receptor antagonist is approved by the FDA for use in anesthesia reversal and benzodiazepine overdose. Research has shown that flumazenil provides relief for most patients whose CSF contains the unknown "somnogen" that enhances the function of GABAA receptors, making them more susceptible to the sleep-inducing effect of GABA. For one patient, daily administration of flumazenil by sublingual lozenge and topical cream has proven effective for several years. A 2014 case report also showed improvement in IH symptoms after treatment with a continuous subcutaneous flumazenil administration. The patient was treated with a short-term subcutaneous administration through 96-hour continuous low dose (4 mg/day) infusion of flumazenil, followed by a slow-release long term subcutaneous administration through flumazenil implant.
Clarithromycin, an antibiotic approved by the FDA for the treatment of infections, was found to return the function of the GABA system to normal in patients with IH. In the pilot study, clarithromycin improved subjective sleepiness in GABA-related hypersomnia. In 2013, a retrospective review evaluating longer-term clarithromycin use showed efficacy in a large percentage of patients with GABA-related hypersomnia. A 2021 Cochrane study determined that the evidence is inadequate to definitively determine the efficacy of clarithromycin in the management of idiopathic hypersomnia. The American Academy of Sleep Medicine's 2021 clinical practice guidelines conditionally suggested its use, especially for those who don't respond to other therapies.
Transcranial direct-current stimulation (tDCS)
Dr. Ferini-Strambi and his colleagues in Milan, Italy, performed neurologic examinations by applying anodal tDCS by placing one electrode over the left dorsolateral prefrontal cortex, with the cathode over the contralateral orbit over 3 weeks period and found that seven of the eight participants (87.5%) reported improvement in their daytime sleepiness, including for up to two weeks after the end of the study.
Transcranial magnetic stimulation (TMS)
Neural networks that regulate arousal and sleep comprise a bottom-up (from the brainstem to the cortex) pathway and a top-down (corticothalamic) pathway. The bottom-up pathway emerges from the ascending reticular arousal system (ARAS) and activates the cortex via well-characterized thalamic and nonthalamic pathways through cholinergic and aminergic neurotransmission. The bottom-up pathway represents the leverage point for pharmaceutical interventions. It is complemented by a corticothalamic top-down pathway, which appears to be modifiable through noninvasive brain stimulation (NIBS) techniques.
A single case report study indicates that high-frequency repetitive transcranial magnetic stimulation (HF rTMS) over the left dorsolateral prefrontal cortex (DLPFC) might represent an alternative choice for symptom control in narcoleptic patients with cataplexy. rTMS may also exert intrinsic effects on hypersomnia in depressed adolescents.
Mazindol
Mazindol is a stimulant similar to amphetamine that "has been shown to be effective in treating hypersomnia in narcoleptics." However, it is not currently approved in the US.
Selegiline
Selegiline, a monoamine oxidase B (MAO-B) inhibitor, works by slowing the breakdown of certain substances in the brain (mostly dopamine, but also serotonin and norepinephrine). It may also be useful, as it is also a metabolic precursor of amphetamine and exerts most of its therapeutic effects through amphetamine metabolism. It is not commonly prescribed for people with narcolepsy because of the high dosage required and potential for severe side effects.
Atomoxetine
Atomoxetine is a norepinephrine reuptake inhibitor (NRI) which increases wakefulness (generally less strongly than the medications which act on dopamine) and which has been argued to have a "clear use in the therapeutic arsenal against narcolepsy and hypersomnia although undocumented by clinical trials."
Ritanserin
Ritanserin is a serotonin antagonist that has "been shown to improve daytime alertness and subjective sleep quality in patients on their usual narcolepsy medications." It is intended as an adjunct (supplement to another main therapeutic agent), and although it is not available in the US, it is available in Europe.
Antidepressants
Antidepressants seems to have some therapeutic effects as they enhance synaptic levels of norepinephrine and serotonin. Further, different medicines are known to augment the activity of one another like as seen in the case of fluoxetine which augmented the activity of methylphenidate when taken together.
Bupropion, a norepinephrine-dopamine reuptake inhibitor (NDRI), which works by inhibiting the reabsorption of two important brain chemicals – norepinephrine and dopamine, is known to have wake-promoting effects.
Fluoxetine, an antidepressant of the selective serotonin reuptake inhibitor (SSRI) class is also known to have mild stimulating effects. It is also known to augment the activity of methylphenidate.
Caffeine
Caffeine is one of the safer nondopaminergic wake-promoting compounds. It is widely used but "has intolerable side effects at high doses (including cardiovascular), and it is generally not efficient enough for patients with hypersomnia or narcolepsy." Although it is commonly used by people with IH or narcolepsy, many people with these disorders report that it has only limited benefit on their sleepiness.
Melatonin
Melatonin is a hormone that the body produces to help regulate sleep. One small study, which used a dose of 2 mg slow-release melatonin at bedtime, found that 50% of participants had “shortened nocturnal sleep duration, decreased sleep inertia, and relieved daytime sleepiness.” Other studies have shown that melatonin synchronizes the circadian rhythms, and improves the “onset, duration and quality of sleep.”
Levothyroxine
There have been some studies suggesting levothyroxine as a possible treatment for IH, especially for patients with subclinical hypothyroidism. This treatment does carry potential risks (especially for patients without hypothyroidism or subclinical hypothroidism), which include cardiac arrhythmia.
Orexin agonists
Orexin-A has been shown to be strongly wake-promoting in animal models, but it does not cross the blood–brain barrier. Suvorexant, a orexin receptor antagonist, has been developed to limit the natural effects of orexin in patients with insomnia. It is therefore possible that a orexin agonist may be similarly developed for the treatment of hypersomnia.
Levodopa
Levodopa (L-DOPA) is an amino acid and is the precursor to the neurotransmitters dopamine, norepinephrine (noradrenaline), and epinephrine (adrenaline). As per a study of six narcoleptic patients it was found that L-DOPA improved vigilance and performance as evaluated by the AVS and the FCRTT, while the capacity to fall asleep rapidly remained unchanged as evaluated by the MSLT. It raises the hypothesis that dopamine may play a role in the physiopathology of excessive daytime sleepiness of this condition.
Carnitine
Carnitine has also been shown to improve narcolepsy symptoms (including daytime sleepiness) by increasing fatty-acid oxidation. Abnormally low levels of acylcarnitine have been observed in patients with narcolepsy. These same low levels have been associated with primary hypersomnia in general in mouse studies. "Mice with systemic carnitine deficiency exhibit a higher frequency of fragmented wakefulness and rapid eye movement (REM) sleep, and reduced locomotor activity." Administration of acetyl-L-carnitine was shown to improve these symptoms in mice. A subsequent human trial found that narcolepsy patients given L-carnitine spent less total time in daytime sleep than patients who were given placebo.
See also
List of investigational sleep drugs
References
External links
Hypersomnia Foundation
Sleep disorders
Hypersomnolence disorders | Idiopathic hypersomnia | [
"Biology"
] | 4,703 | [
"Behavior",
"Sleep",
"Sleep disorders"
] |
39,127,833 | https://en.wikipedia.org/wiki/Cat%20state | In quantum mechanics, the cat state, named after Schrödinger's cat, refers to a quantum state composed of a superposition of two other states of flagrantly contradictory aspects. Generalizing Schrödinger's thought experiment, any other quantum superposition of two macroscopically distinct states is also referred to as a cat state. A cat state could be of one or more modes or particles, therefore it is not necessarily an entangled state. Such cat states have been experimentally realized in various ways and at various scales.
Oftentimes this superposition is described as the system being at both states at the same time, such as the possibilities that a cat would be alive and dead at the same time. This description, however popular, is not correct, since some experimental results depend on the interference of superposed states. For instance, in the well-known double-slit experiment, superposed states give interference fringes, whereas, had the particle been through both appertures, simple addition of single-hole results would obtain.
Cat states over distinct particles
Concretely, a cat state can refer to the possibility that multiple atoms be in a superposition of all spin up and all spin down, known as a Greenberger–Horne–Zeilinger state (GHZ state), which is highly entangled. Since GHZ states are relatively difficult to produce but easy to verify they are often used as a benchmark for different platforms. Such a state for six atoms was realized by a team led by David Wineland at NIST in 2005 and the largest states have since grown to beyond 20.
Optically, the GHZ state can be realized with several distinct photons in a superposition of all polarized vertically and all polarized horizontally.
These have been experimentally realized by a team led by Pan Jianwei at University of Science and Technology of China, for instance, four-photon entanglement, five-photon entanglement, six-photon entanglement, eight-photon entanglement, and five-photon ten-qubit cat state.
This spin up/down formulation was proposed by David Bohm, who conceived of spin as an observable in a version of thought experiments formulated in the 1935 EPR paradox.
Cat states in single modes
In quantum optics, a cat state is defined as the quantum superposition of two opposite-phase coherent states of a single optical mode (e.g., a quantum superposition of large positive electric field and large negative electric field):
where
and
are coherent states defined in the number (Fock) basis. Notice that if we add the two states together, the resulting cat state only contains even Fock state terms:
As a result of this property, the above cat state is often referred to as an even cat state. Alternatively, we can define an odd cat state as
which only contains odd Fock states:
Even and odd coherent states were first introduced by Dodonov, Malkin, and Man'ko in 1974.
Linear superposition of coherent states
A simple example of a cat state is a linear superposition of coherent states with opposite phases, when each state has the same weight:
The larger the value of α, the lower the overlap between the two macroscopic classical coherent states exp(−2α2), and the better it approaches an ideal cat state. However, the production of cat states with a large mean photon number (= |α|2) is difficult. A typical way to produce approximate cat states is through photon subtraction from a squeezed vacuum state. This method usually is restricted to small values of α, and such states have been referred to as Schrödinger "kitten" states in the literature. A method to generate a larger cat state using homodyne conditioning on a number state splitted by a beam splitter was suggested and experimentally demonstrated with a clear separation between the two Gaussian peaks in the Wigner function. More methods have been proposed to produce larger coherent state superpositions through multiphoton subtraction, through ancilla-assisted subtraction, or through multiple photon catalysis steps. Optical methods to "breed" cat states by entangling two smaller "kitten" states on a beamsplitter and performing a homodyne measurement on one output have also been proposed and experimentally demonstrated. If the two "kittens" each have magnitude then when a probabilistic homodyne measurement on the amplitude-quadrature of one beamsplitter output yields a measurement of , the remaining output state is projected into an enlarged cat state where the magnitude has been increased to
Coherent state superpositions have been proposed for quantum computing by Sanders.
Higher-order cat states
It is also possible to control the phase-space angle between the involved coherent amplitudes so that they are not diametrically opposed. This is distinct from controlling the quantum phase relation between the states. Cat states with 3 and 4 subcomponents have been experimentally realized, e.g., one might have a triangular cat state:
or a triangle superposed with vacuum state:
or a square cat state:
The three-component cat states naturally appear as the low-energy eigenstates of three atoms, trapped above a chiral waveguide.
Decoherence
The quantum superposition in cat states becomes more fragile and susceptible to decoherence, the larger they are. For a given well-separated cat state (), an absorption of is sufficient to convert the cat state to a nearly equal mixture of even and odd cat states. For example, with , i.e., ~100 photons, an absorption of just 1% will convert an even cat state to be 57%/43% even/odd, even though this reduces the coherent amplitude by only 0.5%. In other words, the superposition is effectively ruined after the probable loss of just a single photon.
Cat qubit
Cat states can also be used to encode quantum information in the framework of bosonic codes. The idea of using cat qubits as a bosonic code for quantum information processing traces back to Cochrane et al. Quantum teleportation using cat states was suggested by Enk and Hirota and Jeong et al. in view of traveling light fields. Jeong et al. showed that one can discriminate between all of the four Bell states in the cat-state basis using a beam splitter and two photon-number parity detectors, while this task is known to be highly difficult using other optical approaches with discrete-variable qubits. The Bell-state measurement scheme using the cat-state basis and its variants have been found to be useful for quantum computing and communication. Jeong and Kim and Ralph et al. suggested universal quantum computing schemes using cat qubits, and it was shown that this type of approach can be made fault-tolerant.
Bosonic codes
In quantum information theory, bosonic codes encode information in the infinite-dimensional Hilbert space of a single mode.
This is in stark contrast with most encodings for which a 2-dimensional system - a qubit - is used to encode information. The numerous dimensions enable a first degree of redundancy and hence of error protection within a single physical degree of freedom which may consist of the propagating mode of an optical set-up, the vibration mode of a trapped ion or the stationary mode of a microwave resonator. Moreover, the dominant decoherence channel is photon loss and no extra decay channels are known to be added if the number of photons is increased. Hence, to identify a potential error, one needs to measure a single error syndrome, thereby allowing one to realize a significant hardware economy. In these respects, bosonic codes are a hardware efficient path towards quantum error correction.
All the bosonic encodings require non-linearities to be generated, stabilized and measured. In particular, they can't be generated or stabilized with only a linear modes and linear displacements. In practice, ancillary systems are needed for stabilization and error tracking. However, the ancillary systems also have errors, which can in reverse ruin the quantum information. Being immune to these errors is called ‘’fault tolerance’’ and is critical. In particular, even though a linear memory is only subject to photon loss errors, it also experiences dephasing once coupled to a non-linear ancillary system.
Cat codes
Bosonic codes draw their error protection from encoding quantum information in distant locations of the mode phase space. Among these bosonic codes, Schrödinger cat codes encode information as a superposition of coherent states where is the complex amplitude of the field, which are quasi-classical states of the mode.
For instance, the two-component cat code may be defined as:
The computational basis states , and , converge towards the coherent states and when is large.
Another example is the four-component cat code defined as:
Other cat states encoding exist such as squeezed cat codes or pair cat codes in 2-mode system.
2-component cat code
The two basis states of this code and are the coherent states and to a very good approximation when is large. In the language of quantum information processing, cat-state decoherence, mostly originating from single photon loss, is associated with phase-flips. On the contrary, bit-flips bear a clear classical analogue: the random switch between the two coherent states.
Contrary to the other bosonic codes that aim at delocalizing information in both direct space and in reciprocal space, the 2-component cat encoding relaxes one constraint by only delocalizing in one space. The resulting qubit is only protected against one of the two error channels (bit-flips) but consequently the acquired protection is more efficient in terms of required photon number. In order to correct against the remaining error channel (phase-flips), one needs to concatenate with another code in a bias preserving way, such as with a repetition code or a surface code.
As stated above, even though a resonator alone typically suffer only from single photon loss, a finite temperature environment causes single photon gain and the coupling to the non-linear resources effectively induces dephasing. Moreover, single photon losses do not only flip the parity of the cat state but also cause a deterministic decrease of the amplitude of coherent states, the cat “shrinks”. All these effects tend to cause bit-flips. Hence, to protect the encoded states several stabilization procedures were proposed:
dissipative: use engineered dissipation such that its steady states form the cat-qubit manifold.
hamiltonian: use an engineered Hamiltonian such that its degenerate ground states form the cat-qubit manifold
gate-based: regularly re-inflate the cat using optimal control, computer-generated pulses.
The two first approaches are called autonomous since they don't requires active correction, and can be combined. So far, autonomous correction has been proven more fault-tolerant than gate-based correction because of the type of interaction used in gate-based correction.
Bit flip suppression with was demonstrated for two-legged cats with dissipative stabilization at the mere cost of linear increase of phase flip due to single photon loss.
4-component cat code
In order to add first order protection against phase-flips within a single degree of freedom, a higher dimension manifold is required. The 4-component cat code uses the even-parity submanifold of the superposition of 4 coherent states to encode information. The odd-parity submanifold is also 2-dimensional and serves as an error space since a single photon loss switches the parity of the state. Hence, monitoring the parity is sufficient to detect errors caused by single photon loss. Just as in the 2-component cat code, one needs to stabilize the code in order to prevent bit-flips. The same strategies can be used but are challenging to implement experimentally because higher order non-linearities are required.
References
Quantum information science
Quantum states
Fictional cats
Schrödinger's cat | Cat state | [
"Physics"
] | 2,469 | [
"Quantum states",
"Quantum mechanics"
] |
39,129,580 | https://en.wikipedia.org/wiki/Vaimos | Vaimos (Voilier Autonome Instrumenté pour Mesures Océanographiques de Surface) is an autonomous sailing boat with embedded instrumentation for ocean surface measurements.
Its goal is to collect measurements at the surface of the ocean. This robot is the result of a collaboration between ENSTA Bretagne and IFREMER. ENSTA-Bretagne (OSM Team) develops control algorithms and the software architecture, IFREMER (LPO+RDT) builds the mechanics, the embedded instrumentation.
Brest-Douarnenez. One of the longest trips the robot has done is Brest-Douarnenez where Vaimos has done more than 100 km in an autonomous mode. Since Vaimos has made trips more than 350 km long.
In 2013, Vaimos participated in the World Robotic Sailing Championship (WRSC) 2013 in Brest, France
.
Control
The robot follows a desired trajectory which is a sequence of lines. When following a line, the robot has two modes.
the direct mode: the direction to be followed is consistent with the direction of the wind and the distance between the robot and its line corresponds to the GPS error (approx. 10 meters),
the tack mode: the direction to be followed is in the no-go zone and the distance to the line is less than 70 meters.
Two nested control loops are implemented in the computer of Vaimos. The first loop corresponds to a low level controller which tunes the rudder and the sail provided in order to have a desired heading. The second loop generates the desired heading to make the line to follow attractive. With this two loop control strategy, we have the guarantee that the robot will always stay inside the required corridor. This is illustrated the spiral experiment
where the robot has to follow a square spiral
. The proof that the robot will always stay inside its corridor can be performed using interval analysis and set inversion
. The proof assumes that the robot obeys to some uncertain state equations which is not always the case in practice.
References
Robotics
Sailboats
Robots
Unmanned surface vehicles | Vaimos | [
"Physics",
"Technology",
"Engineering"
] | 416 | [
"Machines",
"Robots",
"Automation",
"Physical systems",
"Robotics"
] |
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