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1859894 | https://en.wikipedia.org/wiki/Olefin%20metathesis | Olefin metathesis | In organic chemistry, olefin metathesis is an organic reaction that entails the redistribution of fragments of alkenes (olefins) by the scission and regeneration of carbon-carbon double bonds. Because of the relative simplicity of olefin metathesis, it often creates fewer undesired by-products and hazardous wastes than alternative organic reactions. For their elucidation of the reaction mechanism and their discovery of a variety of highly active catalysts, Yves Chauvin, Robert H. Grubbs, and Richard R. Schrock were collectively awarded the 2005 Nobel Prize in Chemistry.
Catalysts
The reaction requires metal catalysts. Most commercially important processes employ heterogeneous catalysts. The heterogeneous catalysts are often prepared by in-situ activation of a metal halide (MClx) using organoaluminium or organotin compounds, e.g. combining MClx–EtAlCl2. A typical catalyst support is alumina. Commercial catalysts are often based on molybdenum and ruthenium. Well-defined organometallic compounds have mainly been investigated for small-scale reactions or in academic research. The homogeneous catalysts are often classified as Schrock catalysts and Grubbs catalysts. Schrock catalysts feature molybdenum(VI)- and tungsten(VI)-based centers supported by alkoxide and imido ligands.
Grubbs catalysts, on the other hand, are ruthenium(II) carbenoid complexes. Many variations of Grubbs catalysts are known. Some have been modified with a chelating isopropoxybenzylidene ligand to form the related Hoveyda–Grubbs catalyst.
Applications
Olefin metathesis has several industrial applications. Almost all commercial applications employ heterogeneous catalysts using catalysts developed well before the Nobel-Prize winning work on homogeneous complexes. Representative processes include:
The Phillips Triolefin and the Olefin conversion technology. This process interconverts propylene with ethylene and 2-butenes. Rhenium and molybdenum catalysts are used. Nowadays, only the reverse reaction, i.e., the conversion of ethylene and 2-butene to propylene is industrially practiced, however.
Shell higher olefin process (SHOP) produces (alpha-olefins) for conversion to detergents. The process recycles certain olefin fractions using metathesis.
Neohexene production, which involves ethenolysis of isobutene dimers. The catalyst is derived from tungsten trioxide supported on silica and MgO.
1,5-Hexadiene and 1,9-decadiene, useful crosslinking agents and synthetic intermediates, are produced commercially by ethenolysis of 1,5-cyclooctadiene and cyclooctene. The catalyst is derived from Re2O7 on alumina.
Synthesis of pharmaceutical drugs,
Homogeneous catalyst potential
Molecular catalysts have been explored for the preparation of a variety of potential applications. the manufacturing of high-strength materials, the preparation of cancer-targeting nanoparticles, and the conversion of renewable plant-based feedstocks into hair and skin care products.
Types
Some classes of olefin metathesis include:
Cross metathesis (CM)
Ring-opening metathesis (ROM)
Ring-closing metathesis (RCM)
Ring-opening metathesis polymerization (ROMP)
Acyclic diene metathesis (ADMET)
Ethenolysis
Mechanism
Hérisson and Chauvin first proposed the widely accepted mechanism of transition metal alkene metathesis. The direct [2+2] cycloaddition of two alkenes is formally symmetry forbidden and thus has a high activation energy. The Chauvin mechanism involves the [2+2] cycloaddition of an alkene double bond to a transition metal alkylidene to form a metallacyclobutane intermediate. The metallacyclobutane produced can then cycloeliminate to give either the original species or a new alkene and alkylidene. Interaction with the d-orbitals on the metal catalyst lowers the activation energy enough that the reaction can proceed rapidly at modest temperatures.
Olefin metathesis involves little change in enthalpy for unstrained alkenes. Product distributions are determined instead by le Chatelier's Principle, i.e. entropy.
Cross metathesis and ring-closing metathesis are driven by the entropically favored evolution of ethylene or propylene, which can be removed from the system because they are gases. Because of this CM and RCM reactions often use alpha-olefins. The reverse reaction of CM of two alpha-olefins, ethenolysis, can be favored but requires high pressures of ethylene to increase ethylene concentration in solution. The reverse reaction of RCM, ring-opening metathesis, can likewise be favored by a large excess of an alpha-olefin, often styrene. Ring-opening metathesis usually involves a strained alkene (often a norbornene) and the release of ring strain drives the reaction. Ring-closing metathesis, conversely, usually involves the formation of a five- or six-membered ring, which is enthalpically favorable; although these reactions tend to also evolve ethylene, as previously discussed. RCM has been used to close larger macrocycles, in which case the reaction may be kinetically controlled by running the reaction at high dilutions. The same substrates that undergo RCM can undergo acyclic diene metathesis, with ADMET favored at high concentrations. The Thorpe–Ingold effect may also be exploited to improve both reaction rates and product selectivity.
Cross-metathesis is synthetically equivalent to (and has replaced) a procedure of ozonolysis of an alkene to two ketone fragments followed by the reaction of one of them with a Wittig reagent.
Historical overview
"Olefin metathesis is a child of industry and, as with many catalytic processes, it was discovered by accident."
As part of ongoing work in what would later become known as Ziegler–Natta catalysis Karl Ziegler discovered the conversion of ethylene into 1-butene instead of a saturated long-chain hydrocarbon (see nickel effect).
In 1960 a Du Pont research group polymerized norbornene to polynorbornene using lithium aluminum tetraheptyl and titanium tetrachloride (a patent by this company on this topic dates back to 1955),
a reaction then classified as a so-called coordination polymerization. According to the then proposed reaction mechanism a RTiX titanium intermediate first coordinates to the double bond in a pi complex. The second step then is a concerted SNi reaction breaking a CC bond and forming a new alkylidene-titanium bond; the process then repeats itself with a second monomer:
Only much later the polynorbornene was going to be produced through ring opening metathesis polymerisation. The DuPont work was led by Herbert S. Eleuterio. Giulio Natta in 1964 also observed the formation of an unsaturated polymer when polymerizing cyclopentene with tungsten and molybdenum halides.
In a third development leading up to olefin metathesis, researchers at Phillips Petroleum Company in 1964 described olefin disproportionation with catalysts molybdenum hexacarbonyl, tungsten hexacarbonyl, and molybdenum oxide supported on alumina for example converting propylene to an equal mixture of ethylene and 2-butene for which they proposed a reaction mechanism involving a cyclobutane (they called it a quasicyclobutane) – metal complex:
This particular mechanism is symmetry forbidden based on the Woodward–Hoffmann rules first formulated two years earlier. Cyclobutanes have also never been identified in metathesis reactions, which is another reason why it was quickly abandoned.
Then in 1967 researchers led by Nissim Calderon at the Goodyear Tire and Rubber Company described a novel catalyst system for the metathesis of 2-pentene based on tungsten hexachloride, ethanol, and the organoaluminum compound EtAlMe2. The researchers proposed a name for this reaction type: olefin metathesis. Formerly the reaction had been called "olefin disproportionation."
In this reaction 2-pentene forms a rapid (a matter of seconds) chemical equilibrium with 2-butene and 3-hexene. No double bond migrations are observed; the reaction can be started with the butene and hexene as well and the reaction can be stopped by addition of methanol.
The Goodyear group demonstrated that the reaction of regular 2-butene with its all-deuterated isotopologue yielded C4H4D4 with deuterium evenly distributed. In this way they were able to differentiate between a transalkylidenation mechanism and a transalkylation mechanism (ruled out):
In 1971 Chauvin proposed a four-membered metallacycle intermediate to explain the statistical distribution of products found in certain metathesis reactions. This mechanism is today considered the actual mechanism taking place in olefin metathesis.
Chauvin's experimental evidence was based on the reaction of cyclopentene and 2-pentene with the homogeneous catalyst tungsten(VI) oxytetrachloride and tetrabutyltin:
The three principal products C9, C10 and C11 are found in a 1:2:1 regardless of conversion. The same ratio is found with the higher oligomers. Chauvin also explained how the carbene forms in the first place: by alpha-hydride elimination from a carbon metal single bond. For example, propylene (C3) forms in a reaction of 2-butene (C4) with tungsten hexachloride and tetramethyltin (C1).
In the same year Pettit who synthesised cyclobutadiene a few years earlier independently came up with a competing mechanism. It consisted of a tetramethylene intermediate with sp3 hybridized carbon atoms linked to a central metal atom with multiple three-center two-electron bonds.
Experimental support offered by Pettit for this mechanism was based on an observed reaction inhibition by carbon monoxide in certain metathesis reactions of 4-nonene with a tungsten metal carbonyl
Robert H. Grubbs got involved in metathesis in 1972 and also proposed a metallacycle intermediate but one with four carbon atoms in the ring. The group he worked in reacted 1,4-dilithiobutane with tungsten hexachloride in an attempt to directly produce a cyclomethylenemetallacycle producing an intermediate, which yielded products identical with those produced by the intermediate in the olefin metathesis reaction. This mechanism is pairwise:
In 1973 Grubbs found further evidence for this mechanism by isolating one such metallacycle not with tungsten but with platinum by reaction of the dilithiobutane with cis-bis(triphenylphosphine)dichloroplatinum(II)
In 1975 Katz also arrived at a metallacyclobutane intermediate consistent with the one proposed by Chauvin He reacted a mixture of cyclooctene, 2-butene and 4-octene with a molybdenum catalyst and observed that the unsymmetrical C14 hydrocarbon reaction product is present right from the start at low conversion.
In any of the pairwise mechanisms with olefin pairing as rate-determining step this compound, a secondary reaction product of C12 with C6, would form well after formation of the two primary reaction products C12 and C16.
In 1974 Casey was the first to implement carbenes into the metathesis reaction mechanism:
Grubbs in 1976 provided evidence against his own updated pairwise mechanism:
with a 5-membered cycle in another round of isotope labeling studies in favor of the 4-membered cycle Chauvin mechanism:
In this reaction the ethylene product distribution at low conversion was found to be consistent with the carbene mechanism. On the other hand, Grubbs did not rule out the possibility of a tetramethylene intermediate.
The first practical metathesis system was introduced in 1978 by Tebbe based on the (what later became known as the) Tebbe reagent. In a model reaction isotopically labeled carbon atoms in isobutene and methylenecyclohexane switched places:
The Grubbs group then isolated the proposed metallacyclobutane intermediate in 1980 also with this reagent together with 3-methyl-1-butene:
They isolated a similar compound in the total synthesis of capnellene in 1986:
In that same year the Grubbs group proved that metathesis polymerization of norbornene by Tebbe's reagent is a living polymerization system and a year later Grubbs and Schrock co-published an article describing living polymerization with a tungsten carbene complex While Schrock focussed his research on tungsten and molybdenum catalysts for olefin metathesis, Grubbs started the development of catalysts based on ruthenium, which proved to be less sensitive to oxygen and water and therefore more functional group tolerant.
Grubbs catalysts
In the 1960s and 1970s various groups reported the ring-opening polymerization of norbornene catalyzed by hydrated trichlorides of ruthenium and other late transition metals in polar, protic solvents. This prompted Robert H. Grubbs and coworkers to search for well-defined, functional group tolerant catalysts based on ruthenium. The Grubbs group successfully polymerized the 7-oxo norbornene derivative using ruthenium trichloride, osmium trichloride as well as tungsten alkylidenes. They identified a Ru(II) carbene as an effective metal center and in 1992 published the first well-defined, ruthenium-based olefin metathesis catalyst, (PPh3)2Cl2Ru=CHCH=CPh2:
The corresponding tricyclohexylphosphine complex (PCy3)2Cl2Ru=CHCH=CPh2 was also shown to be active. This work culminated in the now commercially available 1st generation Grubbs catalyst.
Schrock catalysts
Schrock entered the olefin metathesis field in 1979 as an extension of work on tantalum alkylidenes. The initial result was disappointing as reaction of with ethylene yielded only a metallacyclopentane, not metathesis products:
But by tweaking this structure to a (replacing chloride by t-butoxide and a cyclopentadienyl by an organophosphine, metathesis was established with cis-2-pentene. In another development, certain tungsten oxo complexes of the type were also found to be effective.
Schrock alkylidenes for olefin metathesis of the type were commercialized starting in 1990.
The first asymmetric catalyst followed in 1993
With a Schrock catalyst modified with a BINOL ligand in a norbornadiene ROMP leading to highly stereoregular cis, isotactic polymer.
| Physical sciences | Organic reactions | Chemistry |
6207163 | https://en.wikipedia.org/wiki/Seismic%20magnitude%20scales | Seismic magnitude scales | Seismic magnitude scales are used to describe the overall strength or "size" of an earthquake. These are distinguished from seismic intensity scales that categorize the intensity or severity of ground shaking (quaking) caused by an earthquake at a given location. Magnitudes are usually determined from measurements of an earthquake's seismic waves as recorded on a seismogram. Magnitude scales vary based on what aspect of the seismic waves are measured and how they are measured. Different magnitude scales are necessary because of differences in earthquakes, the information available, and the purposes for which the magnitudes are used.
Earthquake magnitude and ground-shaking intensity
The Earth's crust is stressed by tectonic forces. When this stress becomes great enough to rupture the crust, or to overcome the friction that prevents one block of crust from slipping past another, energy is released, some of it in the form of various kinds of seismic waves that cause ground-shaking, or quaking.
Magnitude is an estimate of the relative "size" or strength of an earthquake, and thus its potential for causing ground-shaking. It is "approximately related to the released seismic energy".
Intensity refers to the strength or force of shaking at a given location, and can be related to the peak ground velocity. With an isoseismal map of the observed intensities (see illustration) an earthquake's magnitude can be estimated from both the maximum intensity observed (usually but not always near the epicenter), and from the extent of the area where the earthquake was felt.
The intensity of local ground-shaking depends on several factors besides the magnitude of the earthquake, one of the most important being soil conditions. For instance, thick layers of soft soil (such as fill) can amplify seismic waves, often at a considerable distance from the source, while sedimentary basins will often resonate, increasing the duration of shaking. This is why, in the 1989 Loma Prieta earthquake, the Marina district of San Francisco was one of the most damaged areas, though it was nearly 100 km from the epicenter. Geological structures were also significant, such as where seismic waves passing under the south end of San Francisco Bay reflected off the base of the Earth's crust towards San Francisco and Oakland. A similar effect channeled seismic waves between the other major faults in the area.
Magnitude scales
An earthquake radiates energy in the form of different kinds of seismic waves, whose characteristics reflect the nature of both the rupture and the earth's crust the waves travel through. Determination of an earthquake's magnitude generally involves identifying specific kinds of these waves on a seismogram, and then measuring one or more characteristics of a wave, such as its timing, orientation, amplitude, frequency, or duration. Additional adjustments are made for distance, kind of crust, and the characteristics of the seismograph that recorded the seismogram.
The various magnitude scales represent different ways of deriving magnitude from such information as is available. All magnitude scales retain the logarithmic scale as devised by Charles Richter, and are adjusted so the mid-range approximately correlates with the original "Richter" scale.
Most magnitude scales are based on measurements of only part of an earthquake's seismic wave-train, and therefore are incomplete. This results in systematic underestimation of magnitude in certain cases, a condition called saturation.
Since 2005 the International Association of Seismology and Physics of the Earth's Interior (IASPEI) has standardized the measurement procedures and equations for the principal magnitude scales, , , , and .
"Richter" magnitude scale
The first scale for measuring earthquake magnitudes, developed in 1935 by Charles F. Richter and popularly known as the "Richter" scale, is actually the , label ML or ML. Richter established two features now common to all magnitude scales.
First, the scale is logarithmic, so that each unit represents a ten-fold increase in the amplitude of the seismic waves. As the energy of a wave is proportional to A1.5, where A denotes the amplitude, each unit of magnitude represents a 101.5 ≈ 32-fold increase in the seismic energy (strength) of an earthquake.
Second, Richter arbitrarily defined the zero point of the scale to be where an earthquake at a distance of 100 km makes a maximum horizontal displacement of 0.001 mm (1 μm, or 0.00004 in.) on a seismogram recorded with a Wood-Anderson torsion seismograph. Subsequent magnitude scales are calibrated to be approximately in accord with the original "Richter" (local) scale around magnitude 6.
All "Local" (ML) magnitudes are based on the maximum amplitude of the ground shaking, without distinguishing the different seismic waves. They underestimate the strength:
of distant earthquakes (over ~600 km) because of attenuation of the S waves,
of deep earthquakes because the surface waves are smaller, and
of strong earthquakes (over M ~7) because they do not take into account the duration of shaking.
The original "Richter" scale, developed in the geological context of Southern California and Nevada, was later found to be inaccurate for earthquakes in the central and eastern parts of the North American continent (everywhere east of the Rocky Mountains) because of differences in the continental crust. All these problems prompted the development of other scales.
Most seismological authorities, such as the United States Geological Survey, report earthquake magnitudes above 4.0 as moment magnitude (below), which the press describes as "Richter magnitude".
Other "local" magnitude scales
Richter's original "local" scale has been adapted for other localities. These may be labelled "ML", or with a lowercase "l", either Ml, or Ml. (Not to be confused with the Russian surface-wave MLH scale.) Whether the values are comparable depends on whether the local conditions have been adequately determined and the formula suitably adjusted.
Japan Meteorological Agency magnitude scale
In Japan, for shallow (depth < 60 km) earthquakes within 600 km, the Japanese Meteorological Agency calculates a magnitude labeled MJMA, MJMA, or MJ. (These should not be confused with moment magnitudes JMA calculates, which are labeled Mw(JMA) or M(JMA), nor with the Shindo intensity scale.) JMA magnitudes are based (as typical with local scales) on the maximum amplitude of the ground motion; they agree "rather well" with the seismic moment magnitude in the range of 4.5 to 7.5, but underestimate larger magnitudes.
Body-wave magnitude scales
Body-waves consist of P waves that are the first to arrive (see seismogram), or S waves, or reflections of either. Body-waves travel through rock directly.
mB scale
The original "body-wave magnitude" – mB or mB (uppercase "B") – was developed by and to overcome the distance and magnitude limitations of the scale inherent in the use of surface waves. is based on the P and S waves, measured over a longer period, and does not saturate until around M 8. However, it is not sensitive to events smaller than about M 5.5. Use of as originally defined has been largely abandoned, now replaced by the standardized scale.
mb scale
The mb or mb scale (lowercase "m" and "b") is similar to , but uses only P waves measured in the first few seconds on a specific model of short-period seismograph. It was introduced in the 1960s with the establishment of the World-Wide Standardized Seismograph Network (WWSSN); the short period improves detection of smaller events, and better discriminates between tectonic earthquakes and underground nuclear explosions.
Measurement of has changed several times. As originally defined by mb was based on the maximum amplitude of waves in the first 10 seconds or more. However, the length of the period influences the magnitude obtained. Early USGS/NEIC practice was to measure on the first second (just the first few P waves), but since 1978 they measure the first twenty seconds. The modern practice is to measure short-period scale at less than three seconds, while the broadband scale is measured at periods of up to 30 seconds.
mbLg scale
The regional mbLg scale – also denoted mb_Lg, mbLg, MLg (USGS), Mn, and mN – was developed by for a problem the original ML scale could not handle: all of North America east of the Rocky Mountains. The ML scale was developed in southern California, which lies on blocks of oceanic crust, typically basalt or sedimentary rock, which have been accreted to the continent. East of the Rockies the continent is a craton, a thick and largely stable mass of continental crust that is largely granite, a harder rock with different seismic characteristics. In this area the ML scale gives anomalous results for earthquakes that by other measures seemed equivalent to quakes in California.
Nuttli resolved this by measuring the amplitude of short-period (~ 1 second) Lg waves, a complex form of the Love wave that, although a surface wave, he found provided a result more closely related to the scale than the scale. Lg waves attenuate quickly along any oceanic path, but propagate well through the granitic continental crust, and MbLg is often used in areas of stable continental crust; it is especially useful for detecting underground nuclear explosions.
Surface-wave magnitude scales
Surface waves propagate along the Earth's surface, and are principally either Rayleigh waves or Love waves. For shallow earthquakes the surface waves carry most of the energy of the earthquake, and are the most destructive. Deeper earthquakes, having less interaction with the surface, produce weaker surface waves.
The surface-wave magnitude scale, variously denoted as Ms, MS, and Ms, is based on a procedure developed by Beno Gutenberg in 1942 for measuring shallow earthquakes stronger or more distant than Richter's original scale could handle. Notably, it measured the amplitude of surface waves (which generally produce the largest amplitudes) for a period of "about 20 seconds". The scale approximately agrees with at ~6, then diverges by as much as half a magnitude. A revision by , sometimes labeled MSn, measures only waves of the first second.
A modification – the "Moscow-Prague formula" – was proposed in 1962, and recommended by the IASPEI in 1967; this is the basis of the standardized Ms20 scale (Ms_20, Ms(20)). A "broad-band" variant (Ms_BB, Ms(BB)) measures the largest velocity amplitude in the Rayleigh-wave train for periods up to 60 seconds. The MS7 scale used in China is a variant of Ms calibrated for use with the Chinese-made "type 763" long-period seismograph.
The MLH scale used in some parts of Russia is actually a surface-wave magnitude.
Moment magnitude and energy magnitude scales
Other magnitude scales are based on aspects of seismic waves that only indirectly and incompletely reflect the force of an earthquake, involve other factors, and are generally limited in some respect of magnitude, focal depth, or distance. The moment magnitude scale – Mw or Mw – developed by seismologists Thomas C. Hanks and Hiroo Kanamori, is based on an earthquake's seismic moment, M0, a measure of how much work an earthquake does in sliding one patch of rock past another patch of rock. Seismic moment is measured in newton-meters (N m or N⋅m) in the SI, or dyne-centimeters (dyn⋅cm; ) in the older CGS system. In the simplest case the moment can be calculated knowing only the amount of slip, the area of the surface ruptured or slipped, and a factor for the resistance or friction encountered. These factors can be estimated for an existing fault to determine the magnitude of past earthquakes, or what might be anticipated for the future.
An earthquake's seismic moment can be estimated in various ways, which are the bases of the Mwb, Mwr, Mwc, Mww, Mwp, Mi, and Mwpd scales, all subtypes of the generic Mw scale. See for details.
Seismic moment is considered the most objective measure of an earthquake's "size" in regard of total energy. However, it is based on a simple model of rupture, and on certain simplifying assumptions; it does not account for the fact that the proportion of energy radiated as seismic waves varies among earthquakes.
Much of an earthquake's total energy as measured by is dissipated as friction (resulting in heating of the crust). An earthquake's potential to cause strong ground shaking depends on the comparatively small fraction of energy radiated as seismic waves, and is better measured on the energy magnitude scale, Me. The proportion of total energy radiated as seismic waves varies greatly depending on focal mechanism and tectonic environment; and for very similar earthquakes can differ by as much as 1.4 units.
Despite the usefulness of the scale, it is not generally used due to difficulties in estimating the radiated seismic energy.
Two earthquakes differing greatly in the damage done
In 1997 there were two large earthquakes off the coast of Chile. The magnitude of the first, in July, was estimated at , but was barely felt, and only in three places. In October a quake in nearly the same location, but twice as deep and on a different kind of fault, was felt over a broad area, injured over 300 people, and destroyed or seriously damaged over 10,000 houses. As can be seen in the table below, this disparity of damage done is not reflected in either the moment magnitude () nor the surface-wave magnitude (). Only when the magnitude is measured on the basis of the body-wave () or the seismic energy () is there a difference comparable to the difference in damage.
Rearranged and adapted from Table 1 in . Seen also in .
Energy class (K-class) scale
K (from the Russian word класс, 'class', in the sense of a category) is a measure of earthquake magnitude in the energy class or K-class system, developed in 1955 by Soviet seismologists in the remote Garm (Tajikistan) region of Central Asia; in revised form it is still used for local and regional quakes in many states formerly aligned with the Soviet Union (including Cuba). Based on seismic energy (K = log ES, in Joules), difficulty in implementing it using the technology of the time led to revisions in 1958 and 1960. Adaptation to local conditions has led to various regional K scales, such as KF and KS.
K values are logarithmic, similar to Richter-style magnitudes, but have a different scaling and zero point. K values in the range of 12 to 15 correspond approximately to M 4.5 to 6. M(K), M(K), or possibly MK indicates a magnitude M calculated from an energy class K.
Tsunami magnitude scales
Earthquakes that generate tsunamis generally rupture relatively slowly, delivering more energy at longer periods (lower frequencies) than generally used for measuring magnitudes. Any skew in the spectral distribution can result in larger, or smaller, tsunamis than expected for a nominal magnitude. The tsunami magnitude scale, Mt, is based on a correlation by Katsuyuki Abe of earthquake seismic moment () with the amplitude of tsunami waves as measured by tidal gauges. Originally intended for estimating the magnitude of historic earthquakes where seismic data is lacking but tidal data exist, the correlation can be reversed to predict tidal height from earthquake magnitude. (Not to be confused with the height of a tidal wave, or run-up, which is an intensity effect controlled by local topography.) Under low-noise conditions, tsunami waves as little as 5 cm can be predicted, corresponding to an earthquake of M ~6.5.
Another scale of particular importance for tsunami warnings is the mantle magnitude scale, Mm. This is based on Rayleigh waves that penetrate into the Earth's mantle, and can be determined quickly, and without complete knowledge of other parameters such as the earthquake's depth.
Duration and coda magnitude scales
Md designates various scales that estimate magnitude from the duration or length of some part of the seismic wave-train. This is especially useful for measuring local or regional earthquakes, both powerful earthquakes that might drive the seismometer off-scale (a problem with the analog instruments formerly used) and preventing measurement of the maximum wave amplitude, and weak earthquakes, whose maximum amplitude is not accurately measured. Even for distant earthquakes, measuring the duration of the shaking (as well as the amplitude) provides a better measure of the earthquake's total energy. Measurement of duration is incorporated in some modern scales, such as and .
Mc scales usually measure the duration or amplitude of a part of the seismic wave, the coda. For short distances (less than ~100 km) these can provide a quick estimate of magnitude before the quake's exact location is known.
Macroseismic magnitude scales
Magnitude scales generally are based on instrumental measurement of some aspect of the seismic wave as recorded on a seismogram. Where such records do not exist, magnitudes can be estimated from reports of the macroseismic events such as described by intensity scales.
One approach for doing this (developed by Beno Gutenberg and Charles Richter in 1942) relates the maximum intensity observed (presumably this is over the epicenter), denoted I0 (capital I with a subscripted zero), to the magnitude. It has been recommended that magnitudes calculated on this basis be labeled Mw(I0), but are sometimes labeled with a more generic Mms.
Another approach is to make an isoseismal map showing the area over which a given level of intensity was felt. The size of the "felt area" can also be related to the magnitude (based on the work of and ). While the recommended label for magnitudes derived in this way is M0(An), the more commonly seen label is Mfa. A variant, MLa, adapted to California and Hawaii, derives the local magnitude (ML) from the size of the area affected by a given intensity. MI (upper-case letter "I", distinguished from the lower-case letter in Mi) has been used for moment magnitudes estimated from isoseismal intensities calculated per .
Peak ground velocity (PGV) and peak ground acceleration (PGA) are measures of the force that causes destructive ground shaking. In Japan, a network of strong-motion accelerometers provides PGA data that permits site-specific correlation with different magnitude earthquakes. This correlation can be inverted to estimate the ground shaking at that site due to an earthquake of a given magnitude at a given distance. From this a map showing areas of likely damage can be prepared within minutes of an actual earthquake.
Other magnitude scales
Many earthquake magnitude scales have been developed or proposed, with some never gaining broad acceptance and remaining only as obscure references in historical catalogs of earthquakes. Other scales have been used without a definite name, often referred to as "the method of Smith (1965)" (or similar language), with the authors often revising their method. On top of this, seismological networks vary on how they measure seismograms. Where the details of how a magnitude has been determined are unknown, catalogs will specify the scale as "unknown" (variously Unk, Ukn, or UK). In such cases, the magnitude is considered generic and approximate.
An Mh ("magnitude determined by hand") label has been used where the magnitude is too small or the data too poor (typically from analog equipment) to determine a Local magnitude, or multiple shocks or cultural noise complicates the records. The Southern California Seismic Network uses this "magnitude" where the data fail the quality criteria.
A special case is the Seismicity of the Earth catalog of . Hailed as a milestone as a comprehensive global catalog of earthquakes with uniformly calculated magnitudes, they never published the full details of how they determined those magnitudes. Consequently, while some catalogs identify these magnitudes as MGR, others use UK (meaning "computational method unknown"). Subsequent study found many of the values to be "considerably overestimated". Further study has found that most of the magnitudes "are basically for large shocks shallower than 40 km, but are basically for large shocks at depths of 40–60 km". Gutenberg and Richter also used an italic, non-bold M without subscript – also used as a generic magnitude, and not to be confused with the bold, non-italic M used for moment magnitude – and a "unified magnitude" m (bolding added). While these terms (with various adjustments) were used in scientific articles into the 1970s, they are now only of historical interest. An ordinary (non-italic, non-bold) capital "M" without subscript is often used to refer to magnitude generically, where an exact value or the specific scale used is not important.
| Physical sciences | Seismology | Earth science |
23783987 | https://en.wikipedia.org/wiki/Roentgen%20%28unit%29 | Roentgen (unit) | The roentgen or röntgen (; symbol R) is a legacy unit of measurement for the exposure of X-rays and gamma rays, and is defined as the electric charge freed by such radiation in a specified volume of air divided by the mass of that air (statcoulomb per kilogram).
In 1928, it was adopted as the first international measurement quantity for ionizing radiation to be defined for radiation protection, as it was then the most easily replicated method of measuring air ionization by using ion chambers. It is named after the German physicist Wilhelm Röntgen, who discovered X-rays and was awarded the first Nobel Prize in Physics for the discovery.
However, although this was a major step forward in standardising radiation measurement, the roentgen has the disadvantage that it is only a measure of air ionisation, and not a direct measure of radiation absorption in other materials, such as different forms of human tissue. For instance, one roentgen deposits of absorbed dose in dry air, or in soft tissue. One roentgen of X-rays may deposit anywhere from in bone depending on the beam energy.
As the science of radiation dosimetry developed, it was realised that the ionising effect, and hence tissue damage, was linked to the energy absorbed, not just radiation exposure. Consequently new radiometric units for radiation protection were defined which took this into account. In 1953 the International Commission on Radiation Units and Measurements (ICRU) recommended the rad, equal to 100 erg/g, as the unit of measure of the new radiation quantity absorbed dose. The rad was expressed in coherent cgs units. In 1975 the unit gray was named as the SI unit of absorbed dose. One gray is equal to 1 J/kg (i.e. 100 rad). Additionally, a new quantity, kerma, was defined for air ionisation as the exposure for instrument calibration, and from this the absorbed dose can be calculated using known coefficients for specific target materials. Today, for radiation protection, the modern units, absorbed dose for energy absorption and the equivalent dose (sievert) for stochastic effect, are overwhelmingly used, and the roentgen is rarely used. The International Committee for Weights and Measures (CIPM) has never accepted the use of the roentgen.
The roentgen has been redefined over the years. It was last defined by the U.S.'s National Institute of Standards and Technology (NIST) in 1998 as , with a recommendation that the definition be given in every document where the roentgen is used.
History
The roentgen has its roots in the Villard unit defined in 1908 by the American Roentgen Ray Society as "the quantity of radiation which liberates by ionisation one esu of electricity per cm3 of air under normal conditions of temperature and pressure." Using 1 esu ≈ 3.33564 C and the air density of ~1.293 kg/m3 at 0 °C and 101 kPa, this converts to 2.58 × 10−4 C/kg, which is the modern value given by NIST.
1 × 3.33564 × 10−10 × 1,000,000 ÷ 1.293 = 2.58 × 10−4
This definition was used under different names (e, R, and German unit of radiation) for the next 20 years. In the meantime, the French Roentgen was given a different definition which amounted to 0.444 German R.
ICR definitions
In 1928, the International Congress of Radiology (ICR) defined the roentgen as "the quantity of X-radiation which, when the secondary electrons are fully utilised and the wall effect of the chamber is avoided, produce in 1 cc of atmospheric air at 0 °C and 76 cm of mercury pressure such a degree of conductivity that 1 esu of charge is measured at saturation current." The stated 1 cc of air would have a mass of 1.293 mg at the conditions given, so in 1937 the ICR rewrote this definition in terms of this mass of air instead of volume, temperature and pressure. The 1937 definition was also extended to gamma rays, but later capped at 3 MeV in 1950.
GOST definition
The USSR all-union committee of standards (GOST) had meanwhile adopted a significantly different definition of the roentgen in 1934. GOST standard 7623 defined it as "the physical dose of X-rays which produces charges each of one electrostatic unit in magnitude per cm3 of irradiated volume in air at 0 °C and normal atmospheric pressure when ionization is complete." The distinction of physical dose from dose caused confusion, some of which may have led Cantrill and Parker report that the roentgen had become shorthand for 83 ergs per gram (0.0083 Gy) of tissue. They named this derivative quantity the roentgen equivalent physical (rep) to distinguish it from the ICR roentgen.
ICRP definition
The introduction of the roentgen measurement unit, which relied upon measuring the ionisation of air, replaced earlier less accurate practices that relied on timed exposure, film exposure, or fluorescence. This led the way to setting exposure limits, and the National Council on Radiation Protection and Measurements of the United States established the first formal dose limit in 1931 as 0.1 roentgen per day. The International X-ray and Radium Protection Committee, now known as the International Commission on Radiological Protection (ICRP) soon followed with a limit of 0.2 roentgen per day in 1934. In 1950, the ICRP reduced their recommended limit to 0.3 roentgen per week for whole-body exposure.
The International Commission on Radiation Units and Measurements (ICRU) took over the definition of the roentgen in 1950, defining it as "the quantity of X or γ-radiation such that the associated corpuscular emission per 0.001293 gram of air produces, in air, ions carrying 1 electrostatic unit of quantity of electricity of either sign." The 3 MeV cap was no longer part of the definition, but the degraded usefulness of this unit at high beam energies was mentioned in the accompanying text. In the meantime, the new concept of roentgen equivalent man (rem) had been developed.
Starting in 1957, the ICRP began to publish their recommendations in terms of rem, and the roentgen fell into disuse. The medical imaging community still has a need for ionization measurements, but they gradually converted to using C/kg as legacy equipment was replaced. The ICRU recommended redefining the roentgen to be exactly 2.58 × 10−4 C/kg in 1971.
European Union
In 1971 the European Economic Community, in Directive 71/354/EEC, catalogued the units of measure that could be used "for ... public health ... purposes". The directive included the curie, rad, rem, and roentgen as permissible units, but required that the use of the rad, rem and roentgen be reviewed before 31 December 1977. This document defined the roentgen as exactly 2.58 × 10−4 C/kg, as per the ICRU recommendation. Directive 80/181/EEC, published in December 1979, which replaced directive 71/354/EEC, explicitly catalogued the gray, becquerel, and sievert for this purpose and required that the curie, rad, rem and roentgen be phased out by 31 December 1985.
NIST definition
Today the roentgen is rarely used, and the International Committee for Weights and Measures (CIPM) never accepted the use of the roentgen. From 1977 to 1998, the US NIST's translations of the SI brochure stated that the CIPM temporarily accepted the use of the roentgen (and other radiology units) with SI units since 1969. However, the only related CIPM decision shown in the appendix are with regards to the curie in 1964. The NIST brochures defined the roentgen as 2.58 × 10−4 C/kg, to be employed with exposures of x or γ radiation, but did not state the medium to be ionized. The CIPM's current SI brochure excludes the roentgen from the tables of non-SI units accepted for use with the SI. The US NIST clarified in 1998 that it was providing its own interpretations of the SI system, whereby it accepted the roentgen for use in the US with the SI, while recognizing that the CIPM did not. By then, the limitation to x and γ radiation had been dropped. NIST recommends defining the roentgen in every document where this unit is used. The continued use of the roentgen is strongly discouraged by the NIST.
Development of replacement radiometric quantities
Although a convenient quantity to measure with an air ion chamber, the roentgen had the disadvantage that it was not a direct measure of either the intensity of X-rays or their absorption, but rather was a measurement of the ionising effect of X-rays in a specific circumstance; which was dry air at 0 °C and 1 standard atmosphere of pressure.
Because of this the roentgen had a variable relationship to the amount of energy absorbed dose per unit mass in the target material, as different materials have different absorption characteristics. As the science of radiation dosimetry developed, this was seen as a serious shortcoming.
In 1940, Louis Harold Gray, who had been studying the effect of neutron damage on human tissue, together with William Valentine Mayneord and the radiobiologist John Read, published a paper in which a unit of measure, dubbed the "gram roentgen" (symbol: gr) defined as "that amount of neutron radiation which produces an increment in energy in unit volume of tissue equal to the increment of energy produced in unit volume of water by one roentgen of radiation" was proposed. This unit was found to be equivalent to 88 ergs in air. In 1953 the ICRU recommended the rad, equal to 100 erg/g, as the new unit of measure of absorbed radiation. The rad was expressed in coherent cgs units.
In the late 1950s the General Conference on Weights and Measures (CGPM) invited the ICRU to join other scientific bodies to work with the International Committee for Weights and Measures (CIPM) in the development of a system of units that could be used consistently over many disciplines. This body, initially known as the "Commission for the System of Units", renamed in 1964 as the "Consultative Committee for Units" (CCU), was responsible for overseeing the development of the International System of Units (SI). At the same time it was becoming increasingly obvious that the definition of the roentgen was unsound, and in 1962 it was redefined.
The CCU decided to define the SI unit of absorbed radiation in terms of energy per unit mass, which in MKS units was J/kg. This was confirmed in 1975 by the 15th CGPM, and the unit was named the "gray" in honour of Louis Harold Gray, who had died in 1965. The gray was equal to 100 rad. The definition of the roentgen had had the attraction of being relatively simple to define for photons in air, but the gray is independent of the primary ionizing radiation type, and can be used for both kerma and absorbed dose in a wide range of matter.
When measuring absorbed dose in a human due to external exposure, the SI unit the gray, or the related non-SI rad are used. From these can be developed the dose equivalents to consider biological effects from differing radiation types and target materials. These are equivalent dose, and effective dose for which the SI unit sievert or the non-SI rem are used.
Radiation-related quantities
The following table shows radiation quantities in SI and non-SI units:
| Physical sciences | Exposure | Basics and measurement |
4749757 | https://en.wikipedia.org/wiki/Geminal | Geminal | In chemistry, the descriptor geminal () refers to the relationship between two atoms or functional groups that are attached to the same atom. A geminal diol, for example, is a diol (a molecule that has two alcohol functional groups) attached to the same carbon atom, as in methanediol. Also the shortened prefix gem may be applied to a chemical name to denote this relationship, as in a gem-dibromide for "geminal dibromide".
The concept is important in many branches of chemistry, including synthesis and spectroscopy, because functional groups attached to the same atom often behave differently from when they are separated. Geminal diols, for example, are easily converted to ketones or aldehydes with loss of water.
The related term vicinal refers to the relationship between two functional groups that are attached to adjacent atoms. This relative arrangement of two functional groups can also be described by the descriptors α and β.
1H NMR spectroscopy
In 1H NMR spectroscopy, the coupling of two hydrogen atoms on the same carbon atom is called a geminal coupling. It occurs only when two hydrogen atoms on a methylene group differ stereochemically from each other. The geminal coupling constant is referred to as 2J since the hydrogen atoms couple through two bonds. Depending on the other substituents, the geminal coupling constant takes values between −23 and +42 Hz.
Synthesis
The following example shows the conversion of a cyclohexyl methyl ketone to a gem-dichloride through a reaction with phosphorus pentachloride. This gem-dichloride can then be used to synthesize an alkyne.
| Physical sciences | Concepts_2 | Chemistry |
4752147 | https://en.wikipedia.org/wiki/Isenthalpic%20process | Isenthalpic process | An isenthalpic process or isoenthalpic process is a process that proceeds without any change in enthalpy, H; or specific enthalpy, h.
Overview
If a steady-state, steady-flow process is analysed using a control volume, everything outside the control volume is considered to be the surroundings.
Such a process will be isenthalpic if there is no transfer of heat to or from the surroundings, no work done on or by the surroundings, and no change in the kinetic energy of the fluid. This is a sufficient but not necessary condition for isoenthalpy. The necessary condition for a process to be isoenthalpic is that the sum of each of the terms of the energy balance other than enthalpy (work, heat, changes in kinetic energy, etc.) cancel each other, so that the enthalpy remains unchanged. For a process in which magnetic and electric effects (among others) give negligible contributions, the associated energy balance can be written as
If then it must be that
The throttling process is a good example of an isoenthalpic process in which significant changes in pressure and temperature can occur to the fluid, and yet the net sum the associated terms in the energy balance is null, thus rendering the transformation isoenthalpic. The lifting of a relief (or safety) valve on a pressure vessel is an example of throttling process. The specific enthalpy of the fluid inside the pressure vessel is the same as the specific enthalpy of the fluid as it escapes through the valve. With a knowledge of the specific enthalpy of the fluid and the pressure outside the pressure vessel, it is possible to determine the temperature and speed of the escaping fluid.
In an isenthalpic process:
,
.
Isenthalpic processes on an ideal gas follow isotherms, since .
| Physical sciences | Thermodynamics | Physics |
19609175 | https://en.wikipedia.org/wiki/Arctodus | Arctodus | Arctodus is an extinct genus of short-faced bear that inhabited North America during the Pleistocene (~2.5 Mya until 12,800 years ago). There are two recognized species: the lesser short-faced bear (Arctodus pristinus) and the giant short-faced bear (Arctodus simus). Of these species, A. simus was larger, is known from more complete remains, and is considered one of the best known members of North America's extinct Ice Age megafauna. A. pristinus was largely restricted to the Early Pleistocene of the eastern United States, whereas A. simus had a broader range, with most finds being from the Late Pleistocene of the United States, Mexico and Canada. A. simus evolved from A. pristinus, but both species likely overlapped in the Middle Pleistocene. Both species are relatively rare in the fossil record.
Today considered to be an enormous omnivore, Arctodus simus is believed to be one of the largest known terrestrial carnivorans that has ever existed. However, Arctodus, like other bears, was highly sexually dimorphic. Adult A. simus ranged between , with females clustering at ≤, and males around . The largest males stood at at the shoulder, and up to tall on their rear legs. Studies suggest that Arctodus simus browsed on C3 vegetation and consumed browsing herbivores such as deer, camelids, and tapir. A. simus preferred temperate open woodlands but was an adaptable species, taking advantage of many habitats and feeding opportunities.
Arctodus belongs to the Tremarctinae subfamily of bears, which are endemic to the Americas. Of these short-faced bears, Arctodus was the most widespread in North America. However, the genus was restricted to the Pleistocene. A. pristinus went extinct around 300,000 years ago, with A. simus disappearing ~12,800 years ago in the Late Pleistocene extinctions. The cause behind these extinctions is unclear, but in the case of A. pristinus, this was likely due to climate change and competition with other ursids, such as the black bear and Tremarctos floridanus. A. simus likely went extinct due to ecological collapse disrupting the vegetation and prey it relied on.
Taxonomy
Arctodus was first described by Joseph Leidy in 1854, with finds of A. pristinus from the Ashley Phosphate Beds, South Carolina. The scientific name of the genus, Arctodus, derives from Greek, and means "bear tooth". The first fossils of Arctodus simus were found in the Potter Creek Cave, Shasta County, California, by J. A. Richardson in 1878, and were initially described as Arctotherium simum by Edward Drinker Cope in 1879. Historically, all specimens were grouped together under A. pristinus, until a revision by Björn Kurtén in 1967.
In the 19th and early 20th centuries, specimens of Arctodus were occasionally referred to Arctotherium, and vice versa. However, today neither genera are considered to have overlapped, with the closest point of contact being México, with the giant Arctodus simus in Valsequillo, Puebla, and the smaller Arctotherium wingei in the Yucatán Peninsula. Other early researchers believed Arctodus to be a sister lineage of the agriotheriin Indarctos. Sometimes described as the "American cave bear", Arctodus should not be mistaken for the similarly large Eurasian cave bear (Ursus spelaeus). As an ursine, the Eurasian cave bear last shared a common ancestor with the tremarctine Arctodus circa 13.4 million years ago.
Fossils of Arctodus pristinus can be confused with the similarly sized, partially contemporaneous short-faced bear, Tremarctos floridanus. Arctodus has higher crowned and considerably larger teeth than its relative Tremarctos. A. pristinus can be distinguished by broader and taller molars on average, but as they are often worn, differentiation can be difficult. Moreover, diagnosing isolated A. s (such as femora, scapulae, certain vertebrae, ribs, podials) from brown bears can be challenging, as some large brown bears overlap in dimensions with small Arctodus simus. Beyond standard differences between tremarctine and ursine bears, A. simus has a more anterior protocone & extended enamel ridge forming a shearing blade on the maxillary P4. The molars are also shorter & broader in Arctodus than brown bears.
Evolution
Arctodus belongs to the subfamily Tremarctinae, which appeared in North America during the earliest parts of the late Miocene epoch in the form of Plionarctos. The medium-sized Arctodus pristinus, Tremarctos floridanus and Arctotherium sp. evolved from Plionarctos in the Blancan age of North America. The genetic divergence date for Arctodus is ~5 million years ago, around the Miocene-Pliocene boundary, when tremarctine bears, along with other ursids, experienced an explosive radiation in diversity, as C4 vegetation (grasses) and open habitats dominated. The world experienced a major temperature drop and increased seasonality, and a faunal turnover which extinguished 70–80% of North American genera.
Arctodus pristinus was mostly restricted to the more densely forested thermal enclave in eastern North America. A. pristinus has the greatest concentration of fossils in Florida, with the earliest finds being from the Late Blancan Kissimmee River 6 (2.7 - 2.2 Mya) and Santa Fe River 1 sites. During the early Irvingtonian faunal stage, a western population of A. pristinus evolved into the enormous A. simus, with the earliest confirmed records being at least 780,000 years old from the Irvington type locality in California. Correspondingly, A. simus is most plentiful from western North America, albeit preferring mixed habitat such as temperate open woodlands. Their ranges may have met in the Middle Pleistocene of Kansas, with A. simus migrating east in the Late Pleistocene (around the extinction of A. pristinus). Although both Arctodus species co-inhabited North America for at least half a million years during the Middle Pleistocene (A. pristinus went extinct ~300,000 BP), there is no direct evidence of overlap or competition in the fossil record as of yet, as both species established largely separate ranges.
Irvingtonian age (1,900,000 BP - 250,000 BP) specimens of Arctodus simus are particularly sparse. Finds are mostly from California, with additional remains from Texas, Kansas, Nebraska, and Montana. However, A. simus became a pan-continental species in the Rancholabrean faunal stage (Late Pleistocene), sharing that distinction with the black bear. Despite Arctodus simus large temporal and geographic range, fossil remains are comparatively rare (109 finds as of 2010, in otherwise well-sampled localities).
Description
Size
Arctodus pristinus
Around the size of grizzly bears, A. pristinus specimens closely overlap the size of Tremarctos floridanus, with some males of A. pristinus overlapping in size with the females of A. simus. Floridan A. pristinus individuals were calculated to an average of ~. However, the dimensions of some individuals from Port Kennedy Bone Cave and scalientes suggest that northern and western A. pristinus may have been larger than Floridan A. pristinus, being up to .
Arctodus simus
Some A. simus individuals might have been the largest land-dwelling specimens of Carnivora that ever lived in North America. Standing up on its hind legs, A. simus stood , with a maximum vertical arm reach of . When walking on all fours, A. simus stood high at the shoulder, with the largest males being tall enough to look an adult human in the eye. The average weight of A. simus was ~, with the maximum recorded at .
Studies
In a 2010 study, the mass of six A. simus specimens was estimated; half of the specimens weighed between and , with a mean weight of ~850 kg, suggesting larger specimens were probably more common than previously thought. However, the other specimens were calculated to be less than . The weight range calculated from all examined specimens was between 957 kg and . Hypothetically, the largest individuals of A. simus may have approached , or even . However, a 2006 study argued that the maximum size of Arctodus was ~, based on the largest known skull. Additionally, a 1998 study calculated the average weight of Arctodus specimens from the La Brea Tar Pits at ~, smaller than recovered brown bear remains (~, although these remains postdate Arctodus). A 1999 study by Per Christiansen calculated a mean weight of ~ from seven male A. simus limb bones, suggesting large males weighed between and .
Sexual dimorphism
There is much variation in adult size among specimens- the lack of finds, sexual dimorphism, individual variation and potentially ecomorphs could be augmenting the average size of both species of Arctodus. Size differences between specimens of Arctodus simus (such as skull and long bone dimensions) led Kurtén to suggest a larger northern/central subspecies (A. s. yukonensis) and a southern subspecies (A. s. simus). evolving in the Irvingtonian and Rancholabrean respectively. However, the discovery of a very large southern Arctodus simus in Florida and New Mexico (deep within the supposed range of A. s. simus), & possibly Rancho La Brea, and notably small specimens from the Yukon and Vancouver Island, put doubt on this designation. Perceived ecomorphologies are possibly due to the low number of specimens, and sex-biased sampling.
For example, only one baculum (penis bone) has been recovered from over 100 giant short-faced bear sites in North America, although it may belong to a black bear (Potter Cave). None of the specimens assigned to the larger morph (A. s. yukonensis) is from a cave passage, being usually isolated remains from open sites. Furthermore, over 70% of the smaller specimens (once assigned as the A. s. simus subspecies) are from cave deposits where bacula would likely be found if present, suggesting that mostly female individuals of A. simus were using caves. Therefore, in conjunction with ursid sexual dimorphism (e.g. male spectacled bears are 30% - 40% larger than females), the larger, massive Arctodus individuals are often considered male, particularly older males, with the smaller, more lightly built individuals being females. Sexual dimorphism may also explain A. simus teeth (from multiple individuals at the same site) generally clustering into two sizes.
Anatomy
The two species of Arctodus are differentiated not only by size, but also by the shorter snout, greater prognathism, more robust teeth and longer limbs of A. simus, and the relative proportions of each species' molars and premolars. Arctodus pristinus is distinguished from A. simus smaller, narrower, and less crowded teeth. However, the morphologies of both species are otherwise very similar. As a result, differentiating Arctodus simus from Arctodus pristinus can be difficult, as male individuals of Arctodus pristinus can overlap in size with female individuals of Arctodus simus. Arctodus simus superficially resembled living hyaenids in skull shape and relative lengths of the trunk, back and limbs. The most nearly complete skeleton of A. simus found in the United States was unearthed in Fulton County, Indiana; the original bones are in the Field Museum of Natural History, Chicago.
Skull
Members of the Tremarctinae subfamily of bears appear to have a disproportionately short snout compared with most modern bears, giving them the name "short-faced". Arctodus has also been argued to exhibit a wide and shortened rostrum, potentially giving Arctodus a more felid-like appearance. Matheus suggested that a broad snout could have housed a highly developed olfactory apparatus, or accommodated a larger throat passage to bolt down large food items, akin to spotted hyenas. However, this apparent shortness is an illusion caused by the deep snouts and short nasal bones of tremarctine bears compared with ursine bears; Arctodus has a deeper but not a shorter face than most living bears. This characteristic is also shared by the only living tremarctine bear, the omni-herbivorous spectacled bear. Snout deepness could be variable, as specimens from Huntington Reservoir in Utah, and the Hill-Shuler locality, Texas, were noted as being distinctly "short-faced" in comparison with other Arctodus simus individuals.
The orbits of Arctodus are proportionally small compared to the size of the skull, and somewhat laterally orientated (a characteristic of tremarctine bears), more so than actively predatory carnivorans or even the brown bear, suggesting that stereoscopic vision was not a priority. The optic canal and other sphenoidal openings crowd together more in A. simus than in Ursus. As with Tremarctos ornatus, specimens with a large sagittal crest were likely male, whereas females had a reduced or no sagittal crest. Although there are limited samples, the middle ear bones of A. simus are proportionally larger than modern ursine bears, suggesting the species was particularly attuned to low-frequency sounds.
Morphologically, Arctodus simus exhibits masticular characteristics common to herbivorous bears. This includes cheek teeth with large, blunt surface areas, a deep mandible, and large mandibular muscle attachments (which are rare in carnivorous mammals). As herbivorous carnivorans such as Arctodus lack the gut microbiota to efficiently break down plant matter, these features created a high mechanical advantage of the jaw to break down plant matter via extensive chewing or grinding. Although the low mandibular condyle relative to the tooth row (and therefore potential wide gape) of Arctodus simus has been inferred as an adaptation for carnivory, it is also present in the omni-herbivorous spectacled bear. However, both Arctodus pristinus and Tremarctos floridanus have condyles raised well above the plane of the teeth. The purpose of the highly vaulted calvarium and straight cheek bones of Arctodus simus have been similarly disputed.
A 2009 analysis of the mandibular morphology of tremarctine bears found notable differences between Arctodus pristinus and Arctodus simus, with A. simus specimens possessing a concave jaw, large masseter and temporalis muscles, deeper horizontal ramus and a reduced slicing dentition length when compared to A. pristinus. Instead, Arctodus simus was most similar to Arctotherium angustidens- however, both species of Arctodus and Arctotherium angustidens were still comfortably in the "omnivorous" bear cranio-morphotype.
Dentition
The premolars and first molars of Arctodus pristinus are relatively smaller and more widely spaced than those of Arctodus simus. In A. pristinus, the features of the dentition can be quite variable, particularly the M2 molar. An analysis of the Hunter-Schreger bands from the teeth of A. pristinus and A. simus demonstrated an evolutionary trend towards partially reinforced tooth enamel. This has been convergently evolved with giant pandas, agriotheriin bears, and Hemicyon. The dentition of A. simus has been used as evidence of a predatory lifestyle- in particular the large canines, the high-crowned lower first molar, and the possible carnassial shear with the upper fourth premolar. However, the wearing of the molars to a relatively flat & blunt loph (suitable as a crushing platform as per modern omnivorous bears), small shear facet, and the flattened cusps across age ranges (unlike carnivores, which instead have carnassial shears) disagrees with this hypothesis.
Dentition can be a poor indicator of size in A. simus, as some medium-size individuals have teeth that surpass the size of those with the largest skeletons. Additionally, while A. simus evolved from the smaller A. pristinus, their teeth remained generally the same size. A specimen of A. simus from the Seale Pit of the Hill-Shuler locality, Texas, with only two premolars, crowding of the anterior premolar out of line, and a wider and shorter muzzle, was suggested to be an undescribed form of Arctodus.
Post-cranial
Limbs
Researchers have differing interpretations on the limb morphology of Arctodus. A comprehensive 2010 study concluded that the legs of Arctodus weren't proportionally longer than modern bears would be expected to have, and that bears in general are long-limbed animals, obscured in life by their girth and fur. The study concluded the supposed "long-legged" appearance of the bear is largely an illusion created by the animal's relatively shorter back and torso. In fact, Arctodus probably had an even shorter back than other bears, due the necessary ratio between body length and body mass of the huge bear. However, other researchers argue that the limb bones of Arctodus simus are proportionally longer than those of other bears, leading to a "gracile" appearance. Although longer, the proportions still overlap with Ursus, and the limb bones are stouter than in the large-bodied felids (Panthera). Rather than for running, these elongated limb bones may have evolved for increased locomotor efficiency during prolonged travel. This stiff-legged, swinging gait could have been similar to a polar bear's. Some researchers suggest that proportionally longer limbs may be an adaptation for increased vision over tall ground cover in an open habitat, or were used in tearing and pulling down vegetation.
Researchers also disagree when interpreting the humerus of Arctodus simus. Sorkin argued that the pronation of the forearm and the flexion of the wrist and digits, and more lightly muscled forelimbs, all of which are crucial to grasping a large prey animal with the forepaws, were probably less powerful in Arctodus than in either the brown bear or in Panthera. This is due to a weak medial epicondyle and reduced development of the pronator teres muscle. The forelimb of Arctodus could have been in the early stages of cursorial evolution, being capable of more efficient and high-speed straight-line locomotion (relative to extant bears), and was possibly more adept at pursuing large prey than polar and brown bears. On the other hand, some researchers argue that the epicondyles were still well developed, with this wide range of ulna rotation suggests that forearms of Arctodus were powerful and could subdue large prey. A 2013 examination of Rancho La Brean specimens found that they did not possess distally elongated limbs, which discredited cursoriality. Furthermore, the relatively broad humeral & femoral epicondyles were characteristic of diggers and polar bears, and suggested Arctodus simus could have foraged for roots, tubers and ground squirrels and/or had developed forelimb muscles to immobilze moving prey. The shape of the elbow joint, along with an well-developed medial epicondyle which forms an angle with the condyle, and shallower olecranon fossa, would have given Arctodus a higher degree of forelimb dexterity. Originally evolved to facilitate arboreality, other researchers believe that the terrestrial Arctodus (along with Arctotherium and the giant panda) retained this characteristic to assist in foraging for vegetation.
Paws
The paws (metapodials and phalanges) of Arctodus were characteristically long, slender, and more elongated along the third and fourth digits compared to ursine bears. Arctodus' paws were therefore more symmetrical than ursine bears, whose feet have axes aligned with the most lateral (fifth) digit. Also, the first digit of Arctodus was positioned more closely and parallel to the other four digits (i.e. with straight toes, Arctodus had less lateral splaying).
However this is potentially contradicted by possible Arctodus simus trackways from near Lakeview, Oregon, with strong toe splaying, three centrally aligned & evenly spaced toes at the front, and two almost perpendicular lateral toes (80° from the axis of the foot on either side). The trackways suggest that Arctodus had an oval-shaped, undivided pad on its sole, front paws that were slightly larger than its back paws, possessed long claws, and had its hind foot overstep the forefoot when walking, like modern bears. An additional A. simus paw print measuring long and wide has been recovered from White Sands National Park, New Mexico. Some claw marks attributed to Arctodus simus at Riverbluff Cave (as they were four meters above the floor of the cave) were nearly 20 cm in width.
The presence of a partial false thumb in Arctodus simus is a characteristic shared with Tremarctos floridanus and the spectacled bear, and is possibly an ancestral trait. Absent in ursine bears, the false thumb of the spectacled bear has been suggested to assist in herbivorous food manipulation (such as bromeliads, leaves, berries, tree bark & fruits, cactus fruits & pulp, palm hearts & fronds), or arboreality.
Paleopathology
Beyond carbohydrate-associated dental pathologies present in the genus, extensive pathologies have been preserved on the most nearly complete skeleton of Arctodus. The leading hypothesis suggests the Fulton County Arctodus specimen suffered from a syphilis-like (treponemal) disease, or yaws, based on the various lesions present. The same individual records a pathological growth distorting the right humerus, with abscesses are noted between the molars and on both ulna. Hypotheses include syphilis, osteoarthritis, a fungal infection in addition to long term syphilis, or an infected wound. Several specimens from Fairbanks, Alaska, also exhibit either pathological growths or periodontal disease, along with a healed toe bone from Big Bear Cave, Missouri.
Paleobiology
Locomotion
Paul Matheus proposed that Arctodus simus may have moved in a highly efficient, moderate-speed pacing gait, more specialized than modern bears. His research concluded that the large body size, taller front legs, high shoulders, short and sloping back, and long legs of Arctodus also compounded locomotive efficiency, as these traits swelled the amount of usable elastic strain energy in the tendons, and increased stride length, making Arctodus built more for endurance than for great speed. His calculations suggested that Arctodus likely had a top speed of , and based on hyaenid proportions, would shift from singlefoot locomotion to a pace at , and would begin to gallop at , a fairly high speed. Based on other mammals, the optimal pace speed of Arctodus would have been . For comparison, hyenas cross country ~. This mobility would have facilitated travelling across a large home range, which Mattson suggests may have topped . Swimming has also been presented as a hypothesis for the colonization of Vancouver Island by Arctodus simus.
Maturity
Examinations on a mostly full sized young individual of Arctodus simus from an Ozark cave suggest that Arctodus, like other ursids, reached sexual maturity well before full maturity. Comparisons with black bears suggest the Arctodus specimen was either 4–6 years of age if female, or 6–8 years if the specimen was male. Additionally, wear patterns on the individual's teeth are similar to a 4-6 year old Ursus americanus. Fused sutures, epiphyses, and epiphyseal plates, along with tooth eruption, have been used to determine adulthood in Arctodus.
Genetic diversity
An examination of mitochondrial DNA sequenced from specimens of Arctodus simus from Alaska, Alberta, Ohio and the Yukon suggest an extremely low level of genetic diversity among the 23 individuals studied (≤ 44,000 14C BP), with only seven haplotypes recovered. Genetic diversity was comparable to modern endangered fauna, such as the brown kiwi and African cheetah. Explanations include a genetic bottleneck before 44,000 14C BP, or a low level of genetic diversity being a feature of a species which was primarily solitary, with a large home range and relatively small population size. However, this does not entirely preclude genetic diversity in Arctodus simus, with genetic samples from Chiquihuite Cave, Zacatecas indicating a deep divergence with previously studied specimens of A. simus. Additional specimens from the California Channel Islands and Wyoming have been sequenced, but are unassigned.
Haplotype cladogram
Below is a cladogram exploring the relationships between the mitochondrial haplogroups of Arctodus simus. Other than the specimen from Chiquihuite Cave, all haplotypes form a single clade.
Diet
Herbivory
The fact that Arctodus did not significantly differ in dentition or build from modern bears has led most authors to support the hypothesis that the A. simus was omnivorous, like most modern bears, and would have eaten significant amounts of plant matter. Morphologically, Arctodus simus exhibits masticular and dental characteristics which confirms that short-faced bears such as the spectacled bear and Arctodus were adapted to and actively consumed vegetation. This is affirmed by a lack of dental damage associated with carnivory amongst specimens of Arctodus. Dental pathologies which have been found, such as incisor wear & supragingival dental calculus in a young individual from Missouri, and cavities associated with carbohydrate consumption in individuals from the La Brea Tar Pits & Pellucidar Cave (Vancouver Island), further suggest an omnivorous diet for Arctodus simus. Additional morphological adaptations include dexterous forelimbs and a partial false thumb, which would have assisted in foraging for vegetation, along with the body size of large Arctodus simus (~1000 kg) matching or exceeding the expected upper limitations for a terrestrial carnivore (based on the more restrictive energy base for a carnivorous diet).
While features of Arctodus simus morphology suggest herbivory, their close phylogenetic relationship to the omni-herbivorous spectacled bear presents the possibility that these traits may be an ancestral condition of the group. A browsing diet foraged from the canopies of trees and shrubs could have been difficult with the large and flattened rostrum and incisor build of Arctodus, while evidence of digging adaptations in Arctodus' forelimbs and claws (e.g. for rooting) is mixed. Regardless, gross tooth wear suggests consumption of plant matter in the diet of Arctodus simus. The diet of individuals from La Brea was most similar to the spectacled bear, which consumes tough leafy matter, seeded & pitted fruits and occasional protein. Arctodus' tooth wear remained consistent throughout the Pleistocene in La Brea. This indicated a less generalized diet than modern omni-herbivorous black bears, with none of the dental evidence of hard food consumption (such as carcasses or nuts) found in polar bears, black bears and hyenas. Comparisons with the dental microwear of Ursus speleaus suggest dietary differences between the species, with cave bears consuming tougher vegetation than A. simus. Although some researchers argue that herbivory should be more obvious from the isotope data gathered from northern Arctodus, several Arctodus coprolites from The Mammoth Site in South Dakota and Meander Cave at Ni'iinlii'njik Territorial Park, Yukon contain Juniperus seeds (toxic to black & brown bears).
Opportunistic carnivory
Evidence suggests that Arctodus also consumed meat, as evidenced by elevated nitrogen-15 isotope levels (corresponding to protein consumption) and bone damage on contemporary fauna. Additionally, elevated carbon-13 levels (corresponding to C3 resources) from many localities (Alaska, California, San Luis Potosí, Texas, Vancouver Island, and the Yukon) largely suggest browsers (and browsed vegetation) were the core of A. simus' diet.
Arctodus simus' status as a predator is questioned by its gracility and lack of agility, which could have complicated predation upon adult mega-herbivores, and hindered the chasing down of nimbler prey. Nevertheless, larger (male) Arctodus simus are suggested to have been more carnivorous than females, as very large brown bears may not be able to sustain themselves on a vegetarian diet. Furthermore, the much larger frame of A. simus would have provided an advantage in disputes over carcasses.
Studies establish that Arctodus simus would have had a varied diet across its range, and was outcompeted and/or more herbivorous with increased competition from other predators. The extinction of cursorial, hyper-carnivorous Borophagus and Huracan in the more open western North America left a vacant niche, possibly contributing to the evolution of A. simus (along with changes to the herbivore guild).
Bone damage
The bite marks found on many bones of ground sloths (Northrotheriops texanus) and young proboscideans at Leisey Shell Pit in Florida matched the size of the canine teeth of Arctodus pristinus. It is not known if these bite marks are the result of active predation or scavenging. Additionally, Arctodus pristinus was the most common large predator from Port Kennedy Cave, Pennsylvania, where the majority of mastodon remains were juveniles and likely represent accumulated prey.
Arctodus simus has been found in association with proboscidean remains near Frankstown, Pennsylvania (juvenile mastodon), and at The Mammoth Site, South Dakota (Columbian mammoths). However, questions remain as to whether these finds determine a predatory or scavenging relationship, or whether they were simply preserved at the same deposit (possibly at different times). On the other hand, a woolly mammoth specimen from Saltville, Virginia was likely scavenged on by Arctodus simus, as evidenced by a canine gouge through the calcaneus. Several Columbian mammoth bones from a cave near Huntington Reservoir, Utah also record ursid gnaw marks attributed to Arctodus, with an Arctodus specimen preserved in association with the remains.
Importantly, the canines of Panthera atrox overlap in size with Arctodus simus, complicating the identification of tooth marks. However, this is not to discredit all tooth marks attributed to Arctodus, as damaged bones from near Tanana River, Alaska suggest that Arctodus transported megafaunal longbones back to a cave-like den and chewed on them, at a time when lions had a limited overlap with Arctodus in Beringia. Furthermore, a perforated peccary ilium from Sheriden Cave has also been hypothesized as being scavenged by Arctodus simus. Bone damage on a cranial fragment (and possibly the humerus) of an Arctodus individual on Vancouver Island may have been due to cannibalism.
Beringia
Analysis of bones from Alaska showed high concentrations of nitrogen-15, a nitrogen isotope accumulated most strongly in carnivores. Although few specimens exist, there is currently no evidence of the same carbohydrate-related dental pathologies evident in southern populations of Arctodus simus. Based on this evidence, A. simus was suggested to have been more carnivorous in Beringia than the rest of North America (with a preference for herbivores which consumed C3 vegetation, particularly caribou). Increased carnivory may be due to a lower proportion of competitors and probably a lower availability of carbohydrate-rich food supplies across the year in the far northern latitudes. Survival during the cold season for some northern populations of A. simus could have depended on the regular scavenging of ungulate carcasses, as is the case with Alaskan brown bears. Ultimately, an opportunistic foraging strategy including up to 50% vegetation, and the meat of reindeer, muskox, carrion, and possibly some predators, is consistent with the isotopic data and the conclusions of the ecomorphological studies.
Carbon isotope studies
Although elevated nitrogen-15 levels have been argued to indicate carnivory, even the isotope data of the most carnivorous Beringian Arctodus overlapped with modern, typically omni-herbivorous brown bears from Europe, eastern Wyoming, and central Montana, demonstrating that isotope data cannot distinguish between hypercarnivores and omnivores which eat a significant amount of animal matter. Studies are also complicated by a lack of compound-specific data, and isotope data being variable in carbon-13, and nitrogen-15 (due to individual/evolving prey and plant choices, the isotopic composition of the local environment, and nutritional stress). Carbon-13 levels in Arctodus simus (enriched by both plants and prey matter) consistently reflect a diet based on C3 resources, typically found in closed to mixed habitats with at least some tree cover (such as open woodlands). This includes C3 vegetation (leaves, stems, fruits, bark, and flowers from trees, shrubs, and cool season grasses) and the browsers that fed on them, such as deer, camelids, tapir, bison and ground sloths.
{| class="wikitable"
|+Studies
!Location
!Age
!Carbon 13 (δ13C)
!Nitrogen 15 (δ15N)
!Results
|-
|Irvington, California
|Early
Pleistocene
|−14.5
|N/A
|Arctodus simus carbon isotope values from Irvington (along with Fairmead Landfill and McKittrick Tar Pits) are consistent with diet based on C3 resources.
|-
|Fairmead Landfill, California
|Middle Pleistocene
|−11.9
|N/A
|Initially proposed to consume Columbian mammoth, and large ungulates, a 2015 study recalculated Arctodus' carbon isotope values to be closest to C3 vegetation consuming deer and mastodon.
|-
|Cedral, San Luis Potosí
|Late Pleistocene
|−11.8
|N/A
|This Arctodus individual had the strongest δ13C value of its local fauna. Arctodus' carbon isotope value was closest to values from the tapir and Hemiauchenia.
|-
|Natural Trap Cave, Wyoming
|Late Pleistocene
| -13.1
|N/A
|The Natural Trap specimens have the lowest δ13C of the Pleistocene fauna. Comparisons with contemporary Ursus suggest that the seasonality & individual choices within omnivorous diets could result in extreme isotope data, with each tooth not fully reflecting the total isotopic range consumed.
|-
|Channel Islands, California
|Late Pleistocene
| -17.9
|13.2
|Nitrogen isotope signatures suggested a ~19% consumption rate of seals (along with bison and camels). Fossil was likely transported post-mortem from the mainland; a partial reliance on marine resources has been suggested to be as a result of a competitive carnivore guild on mainland California. The marine signal was in between island foxes and bald eagles, most closely resembling Late Pleistocene California condors.
|-
|McKittrick Tar Pits, California
|Late Pleistocene
| -10.9
|N/A
|This carbon isotope value was closest to deer, similar to the one inferred for the Cedral individual.
|-
|Little Box Elder Cave, Wyoming
|Late Pleistocene
| -14.9
|N/A
|Like Natural Trap Cave, the Little Box Elder Cave specimen had distinctly lower δ13C levels, being only higher than Ursus.
|-
|Friesenhahn Cave, Texas
|Late Pleistocene
| -16.5
|9.7
|The Friesenhahn Cave specimen had a nitrogen-15 sample closest to the omnivorous striped skunk.
|-
|Vancouver Island, British Columbia
|Late Pleistocene
| -18.9
|10.6
|A specimen from Cowichan Head, Vancouver Island, had isotopes suggesting a terrestrial diet at a relatively high trophic level.
|-
|Fairbanks, Alaska
|Late Pleistocene
| -18.0
|8.4
|Nitrogen & carbon isotope data from several specimens suggests that Arctodus specialized on reindeer in central Alaska, both before and during the Last Glacial Maximum.
|-
|Dawson, Yukon
|Late Pleistocene
| -18.5
|9.9
|Arctodus nitrogen-15 levels are higher in the Yukon, suggesting a possibly even higher trophic level than other Arctodus in eastern Beringia. However, this contrast likely reflects subtle differences in the isotopic composition of local plants, & muskox in the region, and possibly fellow predators and their kills, complimenting the consumption of reindeer.
|}
Hibernation
Arctodus pristinus specimens have been found in caves such as Port Kennedy, Pennsylvania (where fossils from as many as 36 individuals have been found), and Cumberland Cave, Maryland, often in association with the black bear. This suggests a close association with the biome.
According to a 2003 study, in karst regions, fossils of Arctodus simus have been recovered almost exclusively from cave sites. In the contiguous United States, that ~38% of all sites are from caves (possibly ~50% in western USA) suggests a close association between this species and cave environments. Metabolic denning (hibernation/torpor) is unclear in Arctodus. Like polar bears, male and unmated female A. simus may have forgone denning, leaving maternal denning by females as the preferred explanation behind the recovery of the small, yet relatively complete individuals recovered from caves. However, to date, there are no records of adults with associated offspring from caves. Regardless, Arctotherium angustidens, a fellow giant short-faced bear, has been recovered from a cave in Argentina with offspring.
At Riverbluff Cave, the most abundant claw marks are from Arctodus simus. They are most abundant at the bear beds and their associated passageways, indicating a close relationship with denning. Numerous "bear" beds often preserve Arctodus simus and both Pleistocene and modern American black bears in association (U.a. amplidens and U. a. americanus)- such deposits have been found in Missouri, Oklahoma and Potter Creek Cave, California. These mixed deposits are assumed to have accumulated over time as individual bears (including Arctodus) died during winter sleep. Furthermore, environmental DNA suggests that Arctodus and black bears shared a cave in Chiquihuite cave, Zacatecas. At Labor-of-Love Cave, Nevada, both American black bears and brown bears have been found in association with Arctodus simus. A study in 1985 noted that sympatry between Arctodus and brown bears preserved in caves is rare, with only Little Box Elder Cave, Wyoming and Fairbanks II, Alaska hosting similar remains.
Paleoecology
Arctodus pristinus
Endemic to the late Blancan faunal stage and Irvingtonian faunal stage, Arctodus pristinus was a relatively large tremarctine bear. Sometimes referred to as the eastern short-faced bear, A. pristinus has been found in Florida, Kansas, Maryland, New Mexico, Pennsylvania, South Carolina, and West Virginia in the US, and Aguascalientes in Mexico.''' Possible remains have also been recovered from Arizona. A. pristinus is particularly well known from Florida, especially from the Leisey Shell Pit. Like A. simus and other tremarctine bears, A. pristinus had adaptations for herbivory, and was likely largely herbivorous itself, although Arctodus has been suggested to be generally more carnivorous than contemporary bears.
Eastern North America Arctodus pristinus is considered a biochronological indicator for the period between the Late Blancan and late Irvingtonian periods of Pleistocene Florida- more fossils of Arctodus pristinus are known from Florida (about 150) than anywhere else. In the Early Pleistocene of Blancan Florida, the Santa Fe River 1 site (~2.2 Ma), which Arctodus pristinus inhabited, was a fairly open grassland environment dotted with karst sinks and springs and dominated by longleaf pine flatwoods. Arctodus pristinus co-existed with terror birds, sabertooth cats, giant sloths (Eremotherium, Megalonyx, Paramylodon), giant armadillos (Glyptotherium, Holmesina, Pachyarmatherium), gomphotheres, hyenas, canids (Borophagus, Canis lepophagus), peccaries, llamas, dwarf pronghorns, and three-toed horses. Smaller fauna included condors, rails, ducks, porcupines, and alligators.
Arctodus simus
Evolving from the smaller A. pristinus in the early Irvingtonian faunal stage, scholars today mostly conclude that Arctodus simus was a colossal, opportunistic omnivore, with a flexible, locally adapted diet akin to the brown bear. If Arctodus simus wasn't largely herbivorous, the scavenging of megaherbivore carcasses, and the occasional predatory kill would have complimented the large amounts of vegetation consumed when available.
Sometimes referred to as the bulldog bear, or great short-faced bear, Arctodus simus has been recovered from a comparatively small number of finds in relation to other large carnivorans, with the species suggested to have lived in low population densities. Matheus argues that unlike other Nearctic carnivorans, A. simus did not appear to have an ecological equivalent ("super-huge bear") in the Palearctic realm.Arctodus simus was initially restricted to the western United States during the Irvingtonian. However, in the Rancholabrean faunal stage, A. simus expanded its range from southern Canada to central Mexico in the west, and to Pennsylvania and Florida in the east. A. simus also inhabited eastern Beringia at times, with finds today spanning from northern Alaska to the Yukon. Based on the wide distribution of the species, Arctodus simus inhabited a diversity of climatic conditions and environments. A 2009 study examining megafaunal extinctions in Northern America noted 12 records (<40,000 BP) of Arctodus simus from the Intermontane Plateaus, 7 from the Pacific Mountain System, 6 each from the Interior Plains and Interior Highlands, 3 each from the Atlantic Plains and Rocky Mountain System, and 1 from the Appalachian Highlands.A. simus was relatively plentiful in western North America, with over 50% of specimens from the western contiguous United States (<40,000 BP). Arctodus simus was integral to what has been referred to as the Camelops fauna, or alternatively Camelops/"Navahoceros" fauna, a faunal province centered in western North America. The Camelops fauna was also characterized by shrub-ox, prairie dogs, dwarf pronghorns, Shasta ground sloths, and American lions. The diverse flora of the Camelops faunal province included montane conifers and oak parklands, shrub and grassland that stretched across the North American Cordillera south of Canada, to the Valley of Mexico. This faunal province supported a variety of large grazing and browsing mammals.
Western Mountains
The Pacific Mountain System seems to represent a cradle of evolution for Arctodus simus. The earliest confirmed finds of Arctodus simus are from Irvington, California, which are at least 780,000 years old, but may be older than 1.2Mya. Other Irvingtonian age sites come from California, such as Elsinore, Fairmead, and Murrieta. Older yet disputed remains come from El Casco (1.4Mya).
Despite the shift to aridified, mixed C3-C4 habitats between the Early and Late Pleistocene of the Central Valley (~1Mya to ~15,000 BP), Arctodus simus remained consistent with the consumption of C3 resources. Dire wolves and Arctodus simus were ever present members of the local predator guild throughout the Pleistocene, whereas jaguars, Homotherium, Miracinonyx and Smilodon (Fairmead & Irvington) transitioned to Panthera atrox and coyotes (McKittrick Tar Pits). Although Arctodus could have hunted other closed habitat browsers such as deer (Cervus & Odocoileus), camelids (Hemiauchenia & Camelops), Paramylodon, and peccaries, specimens collected from the La Brea Tar Pits suggest A. simus preferred a herbivorous diet. A. simus is particularly famous from fossils found in the La Brea Tar Pits, with 33 individuals recovered (the most of any locality). As only one juvenile has been found from La Brea, A. simus is suggested to have been solitary. Many more finds come from across California, Vancouver Island, and Washington, where the semi-arid woodland/scrub transitioned to forest-steppe, and open grasslands/heath.
Comparatively, the Rocky Mountain System had the fewest number of specimens of Arctodus simus in western North America. However, one of the youngest dated Arctodus simus is from a cave near Huntington Reservoir, Utah, which sits at an elevation of 2,740m (~9,000 ft). The central and southern Rocky Mountains may have acted as refugia for boreal parkland megafauna from the plateau such as Arctodus simus, with the Huntington specimen being the only confirmed extinct megafauna dated to the Younger Dryas of the Great Basin. Other remains have been found from Wyoming (such as Natural Trap Cave), and Montana.
Intermontane Plateaus
The Intermontane Plateaus had the highest number of Arctodus simus specimens south of the ice sheets. The region has yielded some of the largest specimens of A. simus, including what was once the largest specimen on record, from Salt Lake Valley, Utah. Disputed Irvingtonian remains from eastern California (Victorville and Vallecito Creek) may be as old as 2Mya.
In contrast with other parts of North America, the plateaus received more rainfall during the Late Pleistocene, as glacially cooled air collided with hot desert air. As a result, this greatly expanded the range of subalpine parkland, piñon-juniper & ponderosa woodlands, sagebrush grasslands and pluvial lakes where desert exists today. The mid-Wisconsian U-Bar Cave, New Mexico, was vegetated by sagebrush, grasses, and woodlands. Notable fauna which lived alongside Arctodus simus included Shasta ground sloth, shrub-ox, pronghorns (Stockoceros, Capromeryx), Camelops, Odocoileus, horses, Lynx, puma, black bear, mountain goats, prairie dogs, and Stock's vampire bat. Dire wolves were also found in association with Arctodus simus, and both species are the most common large carnivorans of Rancholabrean New Mexico. Beyond Utah and New Mexico, other important US specimens have also been found in Arizona, eastern California, Idaho, Nevada, and eastern Oregon.
The Intermontane Plateaus extended into central Mexico, with the Mexican Plateau sharing the Late Pleistocene mesic savanna and piñon–juniper woodland ecoregion with the southwestern USA. While Arctodus was limited to the Mexican plateau, the typical tropical thorn scrub and scrub woodland of the plateau was seemingly prime habitat for tremarctine bears. An Arctodus simus individual from Cedral, San Luis Potosí, inhabited closed vegetation, based on the individual's δ13C signature. Consuming C3 resources, its diet may have incorporated local C3 specialists such as tapir, llamas, camels, and Shasta ground sloth along with browsed vegetation. The site, incorporating trees, herbs and cacti, hosted an open gallery forest near grassland or scrub with a humid climate. Similar highland remains have been recovered from Jalisco, Michoacán, Puebla, State of Mexico, and Zacatecas.
Interior Plains
The Interior Plains were composed of temperate steppe grassland, and among the specimens yielded from this region is one of the largest Arctodus simus currently on record, from the banks of the Kansas river. The late Irvingtonian Doeden gravel pits in Montana preserves an open grassland habitat, with riparian woodlands, and likely some shrublands. Arctodus simus co-existed with ground sloths (Megalonyx, Paramylodon), Pacific mastodon, camels, and Bootherium. As bison were yet to migrate into North America, Columbian mammoths and horses dominated these early Illinoian grasslands. Additional Irvingtonian remains are from Kansas, Nebraska and Texas.In the Rancholabrean age, Arctodus simus, grey wolves and coyotes were part of a predator guild throughout the great plains, and were joined by Columbian mammoths, camels, Hemiauchenia, and American pronghorns. While the northern plains aridified into cold steppe (e.g. Mammoth site, South Dakota), the southern plains were a parkland with riparian hackberry forests, and large expanses of mixed grass prairie grasslands grading into wet meadows, with limited seasonality. In the south (Lubbock Lake, Texas), this fauna was joined by Smilodon, dire wolves, grey fox and red fox, preying upon prairie dogs, horses (Equus & Haringtonhippus), peccaries, Odocoileus, Capromeryx, Bison antiquus and Holmesina. Beyond Texas, Arctodus has also been found in Iowa, Kansas, Nebraska, and southern Canada (Alberta & Saskatchewan), which when unglaciated, would have formed a tundra ecosystem with an ice-free corridor to Beringia.
In the lowlands of the eastern Interior plains, the plains transitioned to closed habitat. At the terminal Pleistocene Sheriden Cave, Ohio, a mosaic habitat consisting of marsh, open woodland, and patchy grassland was home to Arctodus simus, Cervalces scotti, caribou, peccaries (Platygonus, Mylohyus), giant beaver, porcupine, and American pine marten. Similar remains have been found in Indiana, and Kentucky.
Interior Highlands
To the south, the Interior Highlands had a very high density of Arctodus simus specimens (second only to the black bear), due to the high rate of preservation in the cave-rich region. Sympatry between the two species is most apparent in Missouri- Arctodus simus has been found in association with black bears at Riverbluff, Bat and Big Bear caves. Big Bear Cave preserves fossilized hair associated with Arctodus. During the Last Glacial Maximum, both bears were joined by dire wolves, coyotes, jaguars, snowshoe hare, groundhogs and beavers at Bat Cave, which also records thousands of Platygonus remains. These fauna inhabited well-watered forest-grassland ecotone with a strong taiga influence, although the region did occasionally cycle through drier, grassier periods. These open woodlands were dominated by pines and spruce, and to a lesser extent by oaks. Additional finds have been recovered from Oklahoma.
Eastern USA
Compared to other regions, Arctodus simus was relatively rare in eastern North America. To the north, the Appalachian Highlands were dominated by taiga. Post-LGM Saltville, Virginia, was a mosaic of grassy/herb laden open areas intermixed with open canopy boreal woodlands (oaks, pines, spruce, birch, firs) and marshes. Inhabiting in this C3 resource dominated environment were Arctodus simus, mastodon, (southernmost) woolly mammoths, Bootherium, horses, caribou, Megalonyx, dire wolves, beavers, Cervalces, and a variety of warm-adapted reptiles, suggesting a more mesic and less seasonal climate than today. Heavy bone damage on a mammoth carcass by both dire wolves and Arctodus suggests a potentially competitive scavenging relationship Beyond Virginia, additional remains have been found in Pennsylvania.
To the south, the subtropical Atlantic Plains covered a great expanse of lowland, from the open deciduous woodlands of the Atlantic coast, to the semi-arid woodland/scrub of Florida, to the spruce-fir conifer forests and open habitat of the Gulf Coastal Plain. Although scarce, this contrast of habitats highlights the adaptability of Arctodus simus. At the Rainbow River and Lake Rousseau localities in Rancholabrean Florida, three Arctodus simus specimens have been recovered, alongside Smilodon, dire wolves, jaguars, ground sloths (Megalonyx, Paramylodon), llamas (Hemiauchenia, Palaeolama), Vero's tapir, giant beaver, capybara, Holmesina, horses, Bison antiquus, mastodon, Columbian mammoths and Tremarctos floridanus, in a climate similar to today's. Furthermore, the abundance of black bears, and particularly Tremarctos floridanus in Florida, has led to a theorized niche partitioning of ursids in Florida, with Tremarctos floridanus being herbivorous, and black bears and Arctodus simus being omnivorous, with Arctodus being possibly more inclined towards carnivory. Additional finds of south-eastern Arctodus simus are from Alabama, Arkansas, Mississippi, South Carolina, and Texas.
Beringia
Largely isolated by the Cordilleran and Laurentide ice sheets, Beringia is considered ecologically separate to the rest of North America, being largely an extension of the mostly open and treeless Eurasian mammoth steppe. However, the occasional opening of an ice-free corridor, and the migration barrier of the Beringian gap, meant that eastern Beringia (Alaska and the Yukon) supported a unique assemblage of fauna, with many endemic North American fauna flourishing. Currently, all specimens of A. simus in Beringia have been dated to a 27,000 year window (50,000 BP - 23,000 BP) from eastern Beringia, while additional undated remains may be of Sangamonian age. Unlike contemporary Beringian carnivorans, A. simus apparently never inhabited western Beringia (and therefore Asia). The largest known skull of A. simus was recovered from the Yukon, and may represent the largest specimen known.
The North Slope of Alaska <40,000 BP (Ikpikpuk and Titaluk rivers) preserves an upland and floodplain environment, with horses, bison then caribou being the most populous herbivores, and woolly mammoths, muskox, elk and saiga antelope more scarce. Cave lions, bears (Ursus arctos and Arctodus simus), and Beringian wolves made up the megafaunal predator guild. Isotope data implies that caribou and muskox were principal components of the carnivorous portion of Arctodus simus' Arctic diet, suggesting that the warmer, wetter vegetation on the margins of the dry mammoth steppe (similar to the moist acidic tundra vegetation which dominates today) was the preferred habitat of Arctodus in Beringia.
Additionally, upon the flooding of the Bering Strait and expansion of moist tundra and peatlands in eastern Beringia during MIS-3, lions, brown bears and Homotherium went regionally extinct ~35,000 BP, whereas wolves and Arctodus persisted. Simultaneously, most megafaunal herbivores in Beringia experienced population bottlenecks, whilst mammoth populations steadily declined. This restriction of prey and habitat could explain the extinctions. However, genetically distinct cave lions and brown bears appear in MIS-2 circa the extinction of Arctodus in a re-emerged Beringia ~23,000 BP, opening up the possibility that some level of competition was at play. The idea that Arctodus had a kleptoparasitic relationship with wolves and Homotherium in Beringia has been explored, with the additional possibility that Arctodus successfully competed against brown bears and Homotherium for access to caribou pre-LGM.
The local extinction of Arctodus in Beringia ~23,000 BP (possibly due to sharp climatic cooling associated with Heinrich Event-2), was much earlier than in other parts of its range. While recolonized by cave lions and brown bears from Eurasia, Arctodus did not repopulate Beringia once the ice-free corridor to the south re-opened later in the Pleistocene.
Map of fossil localities
Relationships with other bears
Arctodus pristinus
In the Early Pleistocene, Arctodus pristinus was much more populous the south-east of North America, whereas the black bear was more common in the north-east. The black bear has inhabited North America since at least the Middle Pleistocene, while Tremarctos floridanus, a tremarctine bear inhabiting western North America at the time, is very similar to A. pristinus in terms of size, skeletal anatomy, and dietary preferences.
Despite this, generally speaking large tremarctine fossils from the Early and Middle Pleistocene of Florida are considered to be A. pristinus, whereas those from the Late Pleistocene of Florida are considered to be T. floridanus. Indeed, black bears and Tremarctos floridanus are believed to have only colonized Florida with the extinction of A. pristinus (both of which only appear in Florida in the Late Pleistocene), however, T. floridanus could yet still be found from older sites in Florida. T. floridanus was possibly an ecological replacement of A. pristinus, with T. floridanus finds being widespread in Rancholabrean Florida and the wider southeastern United States.
Arctodus simus
The most commonly accepted ecological parallel of Arctodus simus in scientific literature is the brown bear. Both brown bears and Arctodus simus exhibit a high degree of dietary variability, and while largely herbivorous, meat can be an important dietary element to certain populations of both species. Additionally, the potential of habitual kleptoparasitism is often noted in Arctodus, with brown bears being opportunistic, curious, and regularly stealing kills from smaller predators. One past theory behind the extinction of Arctodus simus is that A. simus may have been out-competed by brown bears as the latter expanded southwards from eastern Beringia ~13,000 BP, and gradually established itself in North America.
However this has been refuted as new dates establish an extended coexistence, with some isolated A. simus remains being re-evaluated as brown bears. Brown bears (along with lions, bison and red foxes) first emigrated to North America via Beringia during the Illinoian Glaciation, with brown bears first arriving between ~177,000 BP and ~111,000 BP in eastern Beringia. Genetic divergences suggest brown bears first migrated south during MIS-5 (~92,000 - 83,000 BP) upon the opening of the ice-free corridor, with the first fossils being near Edmonton (26,000 BP). On a continent-wide scale, although the brown bear and Arctodus simus were sympatric at times as brown bears spread into North America, Arctodus simus may typically have dominated competitive interactions, and displaced brown bears from specific localities. Additionally, Arctodus' prolonged co-existence with black bears may have put significant constraints on the black bear's evolution.
At the end of the Pleistocene, one reason brown bears persisted where Arctodus simus went extinct was because Arctodus may have been less flexible in adapting to new and rapidly changing environments that impacted the availability or quality of food and habitat. Brown bears and Arctodus have been discovered together in Alaska (then Beringia) between 50,000 BP and 34,000 BP, and in later Pleistocene deposits in Vancouver Island, California, Wyoming and Nevada.
Beringia
Isotope values (δ13C and δ15N) in numerous Beringian Arctodus simus specimens suggests A. simus usually occupied a higher trophic level compared with invading brown bears. While some Beringian brown bears consumed salmon, data from Beringian specimens of Arctodus clustered much more tightly, and suggested that only terrestrial sources of meat were important for Beringian Arctodus. The forcing of a smaller bear into a more herbivorous diet has been compared to the modern relationship between brown bears and American black bears. Where they overlap, black bears take the lower trophic niche, with lower population densities, much smaller territorial ranges, and seasonal migrations. That Arctodus simus (along with local climate change) may have excluded brown bears from eastern Beringia from ~34,000 to ~23,000 BP further suggests that Arctodus may typically have been dominant over brown bears. When Arctodus went extinct in Beringia ~23,000 BP, brown bears recolonized Beringia, but had more carnivorous diets than their Beringian kin pre ~34,000 BP. This bolsters the idea that these bears competed for similar resources and niches. Extinction and repopulation is further evidenced by the high genetic (mitochondrial) diversity of Beringian brown bears in contrast with Beringian Arctodus simus. This contrast in genetic diversity has also been hypothesized to suggest that while female brown bears have a permanent home range, female Arctodus simus may not have (at least not to the same extent).
Vancouver Island
Brown bears, black bears and Arctodus simus all co-existed on Vancouver Island once the island de-glaciated ~14,500 BP. According to an isotope analysis, all three bears relied on terrestrial resources, Arctodus holding an intermediate trophic position between the brown and black bears. This may be an underestimate, as the Arctodus specimens from Vancouver Island are believed to be female; as per brown and black bears, female A. simus may have had a significant decrease in protein consumption compared with male A. simus when co-existing with brown bears. Additionally, an analysis of Arctodus' data suggested that when consuming protein, meat was preferred. While niche-partitioning on Vancouver Island was possible, both Arctodus simus and brown bears appeared to have preferred more open habitats.
Convergent evolution
Both giant short-faced bears Arctodus simus and Arctotherium angustidens reached huge body sizes, in an example of convergent evolution. However, beyond gigantism, there are notable differences between the species. Not only did Arctotherium angustidens reach a higher maximum weight (an exceptional specimen was calculated at ~), A. angustidens was a much more robust animal, in contrast with the gracile Arctodus simus. Excluding the exceptional specimen, Arctotherium angustidens had been calculated to a weight range between and , with the largest specimens of either species being said to be comparable to one another. The panda-relative Agriotherium africanum has also been suggested to share ecomorphological convergences with Arctodus simus. Together with great size, the two species converged on several adaptations, including a skull with a short broad rostrums, premasseteric fossa on the mandible, possible carnassial shears (P4 and m1), and long limbs (relative to body length). These features were also shared by other extinct bears (Agriotherium, Huracan and Arctotherium bonariensis). However, while Agriotherium and Huracan have definitive adaptions for meat-heavy diets stemming from a running, predatory lifestyle, Arctodus simus lacks similar adaptations beyond proportionally longer limbs.
Interactions with humans
One documented interaction with Clovis people is present at the Lubbock Lake Landmark, Texas. A likely already deceased Arctodus simus was processed for subsistence (butchery marks indicated skinning, de-fleshing and disarticulation) and tool production, much in the same way as a mammoth carcass (~13,000 BP / 11,100 14C BP ). Additionally, other remains of the Arctodus simus have been found in association with Paleo-Indian artifacts in Sheriden Cave, Ohio, and Huntington Dam, Utah, with an A. simus footbone fragment from Spalding, Idaho also being charred. The direct relationship between humans and some associated Arctodus remains has been debated. Human hunting and butchery of large megafauna, particularly mammoths and mastodon, would likely have put people in competition with Arctodus simus. Defense against these large bears and the abandonment of carcasses are plausible outcomes, along with the possible caching and disposal of carcass remains underwater to mask its odor from Arctodus.
Migration barrier hypothesis
In the late 1980s, Val Geist hypothesized that "specialist, aggressive, competitive Rancholabrean fauna" such as Arctodus were a barrier for humans (along with other Siberian megafauna such as moose, grey wolves and brown bears) when migrating into North America (both Beringia and below the ice sheets). Male A. simus were the largest and most powerful carnivorous land mammals in North America, with the potential specialization in obtaining and dominating distant and scarce resources. Humans in this hypothesis, though familiar with brown bears, would not have been able to avoid predation or effectively compete with Arctodus simus and other large Pleistocene North American carnivores, making human expansion difficult in Beringia and impossible south of the ice sheets. However, this theory has never been accepted by anthropologists. Paul Matheus argues that there were negligible ecological differences across the mammoth steppe, and that humans successfully competed against and even hunted territorial cave bears, cave hyenas, cave lions, leopards, tigers and wolves in Eurasia before reaching eastern Beringia, making the solitary Arctodus an unlikely impediment to expansion. Indeed, new dates establish an extended co-existence of humans and megafauna such as Arctodus across North America.
Beringia
Humans migrated to North America via the Siberian mammoth steppe, arriving at eastern Beringia (Alaska and the Yukon). However, the migration was halted at the North American Ice Sheet, which separated Beringia and southern North America for most of the Late Pleistocene. Both humans and Arctodus are first dated to ~50,000 BP in Beringia, both from sites in the Yukon, and co-existed until Arctodus went extinct in Beringia ~23,000 BP during the Last Glacial Maximum. This co-existence continued through the regional extinction of other Beringian predators such as cave lions, brown bears and saber-tooth cats. Important sites of pre-LGM human occupation in Beringia include Old Crow Flats and the Klondike, Kuparuk River Valley, and the Bluefish Caves.
Contiguous North America
The human colonization of North America south of the ice sheets further disproves the idea that Arctodus was a migration barrier. The earliest universally accepted pre-Clovis site south of Beringia are the White Sands footprints in New Mexico, dated to ~22,000 cal. BP. Other pre-LGM sites across the Americas, such as Chiquihuite Cave, Valsequillo, El Cedral, Santa Elina, Gault, and Hartley Mammoth Site, affirm that humans proliferated alongside megafauna (such as Arctodus) in southern North America for more than ten thousand years. Humans were definitively widespread across the Americas by at least 15,000 BP.
Extinction
Arctodus pristinus Arctodus pristinus went extinct in the Middle Pleistocene (300,000 years ago), being last recorded from the Coleman 2A site, Florida. The evolution of Arctodus simus, competition with Tremarctos floridanus and black bears, and possibly the transitioning of Pleistocene Florida from a hot, wet, densely forested habitat to a still hot, but drier and much more open biome are thought to be factors behind the gradual disappearance of Arctodus pristinus in the late Irvingtonian faunal stage. There are dubious records of A. pristinus in South Carolina and California from the Late Pleistocene, however these are heavily disputed. Modern research establishes A. pristinus as existing between the Pliocene-Pleistocene boundary and the Middle Pleistocene.
Arctodus simus
With the extinction of Arctodus pristinus, Arctodus simus became the final representative of the genus. Arctodus simus went extinct around 12,800 years ago, and is one of the most recently dated megafauna to go extinct in North America, being reliably dated to within the Pleistocene-Holocene boundary (13,800 BP - 11,400 BP). Both local and regionalized dietary flexibility has been a factor suggested for the species' longevity.
Various factors, including the depletion in number of large herbivores, the diminishing nutritional quality of plants during climate change, and competition with fellow omnivores (humans and brown bears) for food resources, have been suggested as the cause of Arctodus simus' extinction. However, multiple studies put doubt on brown bears being culpable in Arctodus simus' extinction, with the brown bear being more of an ecological replacement that was more adaptable to change. Moreover, there is no systematic evidence that humans hunted large extinct Pleistocene carnivores in North America, and no clear indication of direct human involvement in the extinction of Arctodus simus. Additionally, dental wear evidence from Rancho La Brea does not suggest that food shortages were to blame for the demise large bodied carnivorans such as Arctodus simus.
Climate change
Of the factors discussed, vegetation shifts in the latest Pleistocene may have been particularly unfavorable for Arctodus simus, due to a reduction of quality foraging for subsistence. For example, on Vancouver Island (~13,500 BP), vegetation changed rapidly from open woodlands with abundant lodgepole pine to increasingly closed forests with shade-tolerant spruce, mountain hemlock, and red alder. These changes, effective by ~12,450 BP, point toward cool and moist conditions during the Younger Dryas stadial. Closed forests continued to expand in the early Holocene. Even though Arctodus simus was not restricted to open areas and could inhabit in different environments, the timing of the regional shift from an open pine woodland habitat to a densely forested vegetation implies that these vegetation changes contributed to the local extinction of Arctodus simus, along with many other megafauna.
Low genetic diversity Arctodus simus had a very low level of genetic diversity from most sampled specimens, albeit a sample with a Beringian and temporal bias (<44,000 BP). A loss and/or replacement of mitochondrial DNA lineages before the Last Glacial Maximum, and decrease in population size from a previously genetically diverse population, has been noted in a variety of Eurasian and American Late Pleistocene megafauna. That the individual from Sheriden Cave, Ohio was very closely related to Beringian specimens may further support this idea, as these populations had possibly been isolated from before the Last Glacial Maximum (tens of thousands of years).
A lack of genetic diversity has been attributed to a reduced ability to adapt to environmental conditions. Small population sizes may be characteristic of tremarctine bears- the spectacled bear, while having low levels of genetic diversity, has no signs of a recent genetic bottleneck. However, brown bears had diverse, sympatric source populations in Eurasia, allowing for repopulations/reinvasions into the Americas. If Arctodus simus experienced genetic bottlenecks or local extinctions prior to the Last Glacial Maximum, Arctodus would have been unable to supplement their reduced genetic diversity with new migrants like the brown bear could, making them vulnerable to extinction.
Last dates
The youngest date for Arctodus simus is circa 12,700 BP from Friesenhahn Cave, Texas, calibrated from 10,814 ± 55 radiocarbon years (14C BP). However, this date should be viewed with caution, as analyses suggest the collagen protein was degraded. A vertebra from Bonner Springs, Kansas, was dated to ca. 12,800 BP (based on 10,921 ± 50 radiocarbon years) from well preserved collagen. However, another radiocarbon date from a different laboratory on the same vertebra widens the possible age of the vertebra to between 9,510 and 11,021 14C BP (at 2σ). Nevertheless, a specimen from Huntington Dam, Utah was also dated to ca. 12,800 BP from two radiocarbon dates (10,870 ± 75 & 10,976 ± 40 14C BP) and is therefore considered reliable.
History of research
"Super predator" hypothesis
One past proposal envisaged A. simus as a brutish predator that overwhelmed very large but slow megafauna with its great physical strength. However, despite being very large, its limbs were too gracile for such an attack strategy, significantly more gracile so than Arctotherium angustidens at that.
Due to their long legs, an alternative hypothesis suggested by Björn Kurtén is that it may have hunted by running down Pleistocene herbivores such as wild horses and saiga antelopes, an idea that at one time earned it the name "running bear". However, during pursuit of speedy game animals, the bear's sheer physical mass, inflexible spine and plantigrade gait would be a handicap; modern brown bears can run at the same speed but quickly tire and cannot keep up a chase for long. Correspondingly, although a Arctodus may have been able to reach a maximum speed of , all modern bears have maximum speeds significantly lower than mass-based calculations for speed. As a result, paleontologist Paul Matheus suggests that Arctodus' top speed was . Arctodus skeletons do not articulate in a way that would have allowed for quick turns – an ability required of any predator that survives by chasing down agile prey. Proportionally taller legs, a short trunk, proximally elongated limbs, a stride which had little to no unsupported intervals, small and laterally-orientated eyes, and proportionally short canines ill-suited for spinal and tracheal attacks further complicated ambush hunting as a lifestyle for Arctodus.
Furthermore, the lack of definitive predatory adaptions (such as the absence of laterally compressed canines, and carnassials built for crushing and grinding rather than shearing meat) puts doubt to any species-wide hyper-carnivorous interpretations of Arctodus. The anatomical requirements for a large, cursorial, hyper-carnivorous bear are present in Huracan and Agriotherium, but not Arctodus. Adaptations for predatory behavior are highly divergent in ursids versus other carnivorans, with features such as a short rostrum and long carnassials not being indicative of a predatory lifestyle in Arctodus. Although the only living hyper-carnivorous ursid, the polar bear, also lacks carnassial shears, the species' specialization on small prey and reliance on blubber (rather than coarser flesh) invalidates this comparison with Arctodus. However, both Arctodus simus and polar bears may have had similar overall limb proportions. Regardless, carnivory was likely limited to the regular scavenging of carcasses and opportunistic hunting, as is the case with the modern brown bear.
Specialist kleptoparasite vs Omnivore
The idea that Arctodus was an obligate kleptoparasite was most notably proposed by Paul Matheus. Under this model, A. simus was ill-equipped to be an active predator, having evolved as a specialized scavenger adapted to cover an extremely large home range in order to seek out broadly and unevenly distributed mega-mammal carcasses. There would have been additional selective pressure for increased body size, so that Arctodus could procure and defend carcasses from other large carnivores, some of which were gregarious, or chase them from their kills and steal their food. Matheus calculated that with a hyper-carnivorous diet, a Beringian Arctodus would need to consume ~ of meat per year- the equivalent of 12 bison, 44.6 horses, or 2 woolly mammoths (adjusted for the non-edible portions of the body). Therefore, Arctodus would have had to obtain of flesh/edible carrion every 6.25 days ( per day).
Furthermore, the short rostrum, resulting in increased out-forces of the jaw-closing muscles (temporalis and masseter), may have been an adaptation for cracking bones with their broad carnassials. Such use of the P4 and m1 teeth is supported by the heavy wear on these teeth in old individuals of Arctodus simus and Agriotherium (another giant bear). Additionally, strengthened tooth enamel in Arctodus may have evolved to crack bone. Moreover, at least in Beringia, the conservative growth strategies, long lives and low natural mortality rates of horses and mammoths should have provided somewhat evenly distributed carcasses throughout the year (unlike ruminants such as bison, whose mortality peaks in late winter to early spring). Finally, that Arctodus and the cave hyena did not spread into Siberia and North America respectively suggests some form of competitive exclusion was at play.
Rebuttal
The kleptoparasite hypothesis has been repeatedly challenged. The short, broad rostrum of Arctodus is a characteristic also shared with the sun bear and the spectacled bear, which are both omnivorous. Specialized scavengers like hyenas show distinctive patterns of molar damage from cracking bones. Based on lack of "bone-cracking" wear in specimens from Rancho La Brea, researchers in 2013 concluded that Arctodus simus was not a specialized scavenger. Of living bears, this population of A. simus showed the most similar tooth wear patterns to its closest living relative, the spectacled bear, which can have a highly varied diet ranging from omnivory to almost pure herbivory.
Additionally, severe tooth crown fractures and alveolar infections were found in the South American giant short-faced bear (Arctotherium angustidens). These were interpreted as evidence of feeding on hard materials (e.g. bones), which could tentatively indicate for these bears the regular scavenging of ungulate carcasses obtained through kleptoparasitism. However, such dental pathologies were not observed in various specimens of A. simus, other than the strong wear facets of old individuals. Instead, recovered dental damage (incisor wear, dental calculus & cavities) is herbivorous in origin. Moreover, researchers in 2015 reviewing links between canine breakage, microwear texture patterns and carnivorans from La Brea found that A. simus consumed foods softer yet tougher than black bears and polar bears, avoided hard/brittle foods such as bone, and reaffirmed affinities between A. simus and modern, largely herbivorous spectacled bears. In addition to hyenas, many other fauna did not cross the Rancholabrean Beringian gap, such as the American badger, Bootherium and the woolly rhino).
Furthermore, the relative lack of Arctodus remains at predator traps such as the La Brea Tar Pits, suggests that Arctodus did not regularly compete for carcasses. Although La Brea has produced more Arctodus simus specimens than any other site, Arctodus represents only 1% of all carnivorans in the pits. While more abundant than brown bears and black bears, Arctodus was calculated to its baseline continental abundance, contrasting with the overabundance of other large carnivorans. A similar rate (~0.9%) of relative abundance was calculated for Arctodus compared to other megafauna at the Natural Trap Cave in Wyoming by 1993. Additionally, isotope analyses of Beringian Arctodus specimens suggest that Arctodus had a low consumption rate of horses and mammoths in Beringia, despite those species making up ~50% of the available biomass in Beringia. Further evidence comes from the evolution of brain size relative to body size- bears with high caloric diets and which do not exhibit dormancy showed a weak but significant correlation with bigger relative brain size. Arctodus simus plotted in between the likely hypercarnivorous Cephalogale, and the almost exclusively herbivorous Eurasian cave bear and Indarctos, suggesting omnivory.
| Biology and health sciences | Bears | Animals |
19612238 | https://en.wikipedia.org/wiki/Public%20transport%20bus%20service | Public transport bus service | Public transport bus services are generally based on regular operation of transit buses along a route calling at agreed bus stops according to a published public transport timetable.
History of buses
Origins
While there are indications of experiments with public transport in Paris as early as 1662, there is evidence of a scheduled "bus route" from Market Street in Manchester to Pendleton in Salford UK, started by John Greenwood in 1824.
Another claim for the first public transport system for general use originated in Nantes, France, in 1826. Stanislas Baudry, a retired army officer who had built public baths using the surplus heat from his flour mill on the city's edge, set up a short route between the center of town and his baths. The service started on the Place du Commerce, outside the hat shop of a M. Omnès, who displayed the motto Omnès Omnibus (Latin for "everything for everybody" or "all for all") on his shopfront. When Baudry discovered that passengers were just as interested in getting off at intermediate points as in patronizing his baths, he changed the route's focus. His new voiture omnibus ("carriage for all") combined the functions of the hired hackney carriage with a stagecoach that travelled a predetermined route from inn to inn, carrying passengers and mail. His omnibus had wooden benches that ran down the sides of the vehicle; passengers entered from the rear.
In 1828, Baudry went to Paris, where he founded a company under the name Entreprise générale des omnibus de Paris, while his son Edmond Baudry founded two similar companies in Bordeaux and in Lyon.
A London newspaper reported on July 4, 1829, that "the new vehicle, called the omnibus, commenced running this morning from Paddington to the City", operated by George Shillibeer.
The first omnibus service in New York began in 1829, when Abraham Brower, an entrepreneur who had organized volunteer fire companies, established a route along Broadway starting at Bowling Green. Other American cities soon followed suit: Philadelphia in 1831, Boston in 1835 and Baltimore in 1844. In most cases, the city governments granted a private company—generally a small stableman already in the livery or freight-hauling business—an exclusive franchise to operate public coaches along a specified route. In return, the company agreed to maintain certain minimum levels of service.
In 1832, the New York omnibus had a rival when the first trams, or streetcars started operation along Bowery, which offered the excellent improvement in amenity of riding on smooth iron rails rather than clattering over granite setts, called "Belgian blocks". The streetcars were financed by John Mason, a wealthy banker, and built by an Irish-American contractor, John Stephenson. The Fifth Avenue Coach Company introduced electric buses to Fifth Avenue in New York in 1898.
In 1831, New Yorker Washington Irving remarked of Britain's Reform Act (finally passed in 1832): "The great reform omnibus moves but slowly." Steam buses emerged in the 1830s as competition to the horse-drawn buses.
The omnibus extended the reach of the emerging cities. The walk from the former village of Paddington to the business heart of London in the City was a long one, even for a young man in good condition. The omnibus thus offered the suburbs more access to the inner city. The omnibus encouraged urbanization. Socially, the omnibus put city-dwellers, even if for only half an hour, into previously-unheard-of physical intimacy with strangers, squeezing them together knee-to-knee. Only the very poor remained excluded. A new division in urban society now came to the fore, dividing those who kept carriages from those who did not. The idea of the "carriage trade", the folk who never set foot in the streets, who had goods brought out from the shops for their appraisal, has its origins in the omnibus crush.
Motorbus
John D. Hertz founded the Yellow Coach Manufacturing Company in 1923 and then sold a majority of shares to General Motors in 1925.
From the 1920s, General Motors and others started buying up streetcar systems across the United States with a view to replacing them with buses in what became known as the Great American Streetcar Scandal. This was accompanied by a continuing series of technical improvements: pneumatic "balloon" tires during the early 1920s, monocoque body construction in 1931, automatic transmission in 1936, diesel engines in 1936, 50+ passengers in 1948, and air suspension in 1953.
The arrest of Rosa Parks in 1955 for not giving up her seat to a white man on a public bus is considered one of the catalysts of the Civil Rights Movement within the United States.
Types of services
The names of different types of bus services vary according to local tradition or marketing, although services can be classified into basic types based on route length, frequency, the purpose of use and type of bus used.
Urban transport or regional
Transit bus is the most common type of public transport bus service and is used to transport large numbers of people in urban areas, or to and from the suburbs to population centres. These buses normally run on fixed routes within an urban area.
Park and ride bus services are designed to provide an onward passenger journey from a parking lot. These may be branded as shuttle or express services, or part of the standard bus network.
Share taxi bus services are designed to run as flexible high capacity vehicles usually using minibuses to any point of a person's wish instead of a fixed route. Most common examples of share taxis include public light buses in Hong Kong where the red topped ones act as a share taxi as opposed to green topped ones which are on fixed routes.
Feeder bus services are designed to pick up passengers in a certain locality and take them to a transfer point where they make an onward journey on a trunk service. This can be another bus, or a rail-based service such as a tram, rapid transit or train. Feeder buses may act as part of a wider local network, or a regional coach network.
Bus rapid transit (BRT) is the application of a range of infrastructure and marketing measures to produce public transport bus services that approach the operating characteristics and capacity of rapid transit systems.
Express bus service
An express bus service (also known as express commuter service, commuter bus service, or suburban bus service) is a fixed-route bus service that is intended to run faster than normal bus services between the same two commuter or destination points, typically on longer-distance routes.
Express buses operate on a faster schedule by not making as many stops as normal bus services and often taking quicker routes, such as along freeways, or by using dedicated lanes or roadways. Express buses may also operate out of park and rides, in some cases only during rush hour in the peak direction.
Fares on express bus services may be higher than normal parallel services. Many express buses act as precursors to bus rapid transit lines and employ a proof-of-payment scheme, requiring passengers to purchase tickets before boarding the bus, speeding up the service. These services may also use suburban coaches that feature amenities like comfortable seating and wireless Internet service, particularly on routes that travel long distances at higher speeds without stopping.
In many cases, an express bus service is identified by a letter before or after the regular route number. For example, in Sydney, the letters L (as in L90), E (as in E70) and X (as in 610X or X84). L indicates that the bus runs along the normal route, while E and X indicate that the bus runs along a more direct route. In New York City, express buses operate using coaches from Motor Coach Industries and Prevost Car, and all except the operate along highways, sometimes for a large portion of the route. For example, the Super Expresses, the all operate on highways for most of their route. Many transit systems may also use a specific number before or after the regular route number. For example, in Toronto, the number "9" (as in 995) goes before the regular route number to display an express bus service.
Long distance transport
Long-distance coach services (US: Intercity bus line) are bus services operated over long distances between cities. These services can form the mainstay of the travel network in countries with poor railway infrastructure. Different coach operators may band together on a franchise or connecting basis to offer a branded network that covers large distances, such as Trailways and National Express. These networks can even operate internationally, such as Eurolines of Europe. Interurban bus services are primarily aimed at linking together one or more urban centres, and as such are often run as express services while travelling in the intermediate rural areas, or even only call at two terminal points as a long distance shuttle service. Some interurban services may be operated as high specification luxury services, using coaches, in order to compete with railways, or link areas not rail connected. Interurban services may often terminate in central bus stations rather than on street stops. Other interurban services may specifically call at intermediate villages and may use slower transit buses or dual purpose buses.
Specialist services
School buses transport children to and from school. While many countries and school districts organise their own services, as school buses or charter buses, in some areas school bus services are implemented as special journeys on the normal public timetable, specially timed and routed to arrive and depart in coordination with the school bell. Only the latter is commonly referred to as "public transport".
Shuttle buses are any type of bus service intended primarily to shuttle passengers between two fixed points. These can be bus or coach operated, but are usually short or medium distance journeys taking less than an hour. Shuttle buses will usually link with other transport hubs, such as airport shuttle buses. A common use of a shuttle bus is in towns or cities with multiple terminal train stations or bus stations, for passenger interconnections. "Shuttle" as a brand name is applied variously across several types of service.
Post bus services are services that also carry mail, often on rural routes.
Rail replacement bus services are often chartered by railway companies as alternate means of transport for rail passengers. This can be pre-planned to cover for scheduled track maintenance or other planned closures, or to cover for unplanned closures such as derailments.
Operation
Scheduling
Many public bus services are run to a specific timetable giving specific times of departure and arrival at waypoints along the route. These are often difficult to maintain in the event of traffic congestion, breakdowns, on/off bus incidents, road blockages or bad weather. Predictable effects such as morning and evening rush hour traffic are often accounted for in timetables using the past experience of the effects, although this then prevents the opportunity for drafting a 'clock face' timetable where the time of a bus is predictable at any time through the day. Predictable short term increases in passenger numbers may be dealt with by providing "duplicate" buses, where two or more buses operate the same slot in the timetable. Unpredictable problems resulting in delays and gaps in the timetabled service may be dealt with by 'turning' a bus early before it reaches it terminus, so that it can fill a gap in the opposite direction, meaning any passengers on the turned bus need to disembark and continue on a following bus. Also, depending on the location of the bus depot, replacement buses may be dispatched from the depot to fill in other gaps, starting the timetable part way along the route.
There is a common cliché that people "wait all day, and then three come along at once", in relation to a phenomenon where evenly timetabled bus services can develop a gap in service followed by buses turning up almost simultaneously. This occurs when the rush hour begins and numbers of passengers at a stop increases, increasing the loading time, and thus delay scheduled service. The following bus then catches up because it begins to be delayed less at stops due to fewer passengers waiting. This is called bus bunching. This is prevented in some cities such as Berlin by assigning every stop arrival times where scheduled buses should arrive no earlier than specified.
Some services may have no specific departure times, the timetable giving the frequency of service on a route at particular phases of the day. This may be specified with departure times, but the over-riding factor is ensuring the regularity of buses arriving at stops. These are often the more frequent services, up to the busiest bus rapid transit schemes. For headway-based schemes, problems can be managed by changing speed, delaying at stops and leap-frogging a bus boarding at a stop.
Services may be strictly regulated in terms of level of adherence to timetables, and how often timetables may be changed. Operators and authorities may employ on-street bus inspectors to monitor adherence in real time. Service operators often have a control room, or in the case of large operations, route controllers, who can monitor the level of service on routes and can take remedial action if problems occur. This was made easier with the technological advances of two way radio contact with drivers, and vehicle tracking systems.
Urban land-use planning policies are essential for the success of bus transit systems, particularly as mass transit is not feasible in low-density communities. Transportation planners estimate that to support local bus service every thirty minutes, there must be a residential housing density of seven dwelling units per acre.
Fixed infrastructure
Bus services have led to the implementation of various types of infrastructure now common in many urban and suburban settings. The most prevalent example is the ubiquitous bus stop. Large interchanges have required the building of bus stations. In roads and streets, infrastructure for buses has resulted in modifications to the kerb line such as protrusions and indentations, and even special kerb stones. Entire lanes or roads have been reserved for buses in bus lanes or busways. Bus fleets require large storage premises often located in urban areas, and may also make use of central works facilities.
Management
The level and reliability of bus services are often dependent on the quality of the local road network and levels of traffic congestion, and the population density. Services may be organised on tightly regulated networks with restrictions on when and where services operate, while other services are operated on an ad hoc basis in the model of share taxis.
Increasingly, technology is being used to improve the information provided to bus users, with vehicle tracking technologies to assist with scheduling, and to achieve real-time integration with passenger information systems that display service information at stops, inside buses, and to waiting passengers through personal mobile devices or text messaging.
Fare models
Bus drivers may be required to conduct fare collection, inspect a travel pass or free travel pass, or oversee stored-value card debiting. This may require the fitting of equipment to the bus. Alternatively, this duty and equipment may be delegated to a conductor who rides on the bus. In other areas, public transport buses may operate on a zero-fare basis, or ticket validation may be through the use of on-board/off-board proof-of-payment systems, checked by roving ticket controllers who board and alight buses at random.
In some competitive systems, an incumbent operator may introduce a "low-cost unit" paying lower wages, in order to be able to offer lower fares, using older buses cascaded from a main fleet to also reduce costs. In some sectors, operators such as Megabus (both in the UK and in North America) have attempted to emulate the low-cost airlines model in order to attract passengers through low fares, by offering no-frills bus services.
Ownership
Public transport bus operation is differentiated from other bus operation by the fact the owner or driver of a bus is employed by or contracted to an organisation whose main public duty or commercial interest is to provide a public transport service for passengers to turn up and use, rather than fulfilling private contracts between the bus operator and user. Public transport buses are operated as a common carrier under a contract of carriage between the passenger and the operator.
The owners of public transport buses may be the municipal authority or transit authority that operates them, or they may be owned by individuals or private companies who operate them on behalf of the authorities on a franchise or contract basis. Other buses may be run entirely as private concerns, either on an owner-operator basis, or as multi-national transport groups. Some countries have specifically deregulated their bus services, allowing private operators to provide public bus services. In this case, an authority may make up the shortfall in levels of private service provision by funding or operating 'socially necessary' services, such as early or late services, on the weekends, or less busy routes. Ownership/operation of public transport buses can also take the form of a charitable operation or not for profit social enterprises.
Larger operations may have fleets of thousands of vehicles. At its peak in the 1950s, the London Transport Executive owned a bus fleet of 8,000 buses, the largest in the world. Many small operators have only a few vehicles or a single bus owned by an owner-driver.
Andhra Pradesh State Road Transport Corporation holds the Guinness world record of having the largest fleet of buses with 22,555 buses.
Regulation
In all cases in the developed world, public transport bus services are usually subject to some form of legal control in terms of vehicle safety standards and method of operation, and possibly the level of fares charged and routes operated.
Bus services are being made accessible, often in response to rules and regulations in disability discrimination laws. This has resulted in the introduction of paratransit services and low-floor buses to support passengers who are elderly, have a disability, or a medical condition.
Some transit agencies have also started to install bike racks in the front of buses that usually holds two bicycles. Passengers would be able to place their bicycle on the racks when riding to avoid taking up space during rush hour.
Safety
The research conducted in Montreal (Canada) showed that travelling by bus is safer than travelling by car, for vehicle occupants but also for pedestrians and cyclists. There were 16 times more injured car occupants than bus occupants. Most pedestrians (95%) and cyclists (96%) were injured by a car. Looking at major injuries only (excluding minor injuries), there were 28 times more injured car occupants than bus occupants. Cars were associated with three cyclist deaths and 42 pedestrian deaths while buses were associated with no cyclist deaths and four pedestrian deaths.
| Technology | Motorized road transport | null |
19616384 | https://en.wikipedia.org/wiki/Abstract%20algebra | Abstract algebra | In mathematics, more specifically algebra, abstract algebra or modern algebra is the study of algebraic structures, which are sets with specific operations acting on their elements. Algebraic structures include groups, rings, fields, modules, vector spaces, lattices, and algebras over a field. The term abstract algebra was coined in the early 20th century to distinguish it from older parts of algebra, and more specifically from elementary algebra, the use of variables to represent numbers in computation and reasoning. The abstract perspective on algebra has become so fundamental to advanced mathematics that it is simply called "algebra", while the term "abstract algebra" is seldom used except in pedagogy.
Algebraic structures, with their associated homomorphisms, form mathematical categories. Category theory gives a unified framework to study properties and constructions that are similar for various structures.
Universal algebra is a related subject that studies types of algebraic structures as single objects. For example, the structure of groups is a single object in universal algebra, which is called the variety of groups.
History
Before the nineteenth century, algebra was defined as the study of polynomials. Abstract algebra came into existence during the nineteenth century as more complex problems and solution methods developed. Concrete problems and examples came from number theory, geometry, analysis, and the solutions of algebraic equations. Most theories that are now recognized as parts of abstract algebra started as collections of disparate facts from various branches of mathematics, acquired a common theme that served as a core around which various results were grouped, and finally became unified on a basis of a common set of concepts. This unification occurred in the early decades of the 20th century and resulted in the formal axiomatic definitions of various algebraic structures such as groups, rings, and fields. This historical development is almost the opposite of the treatment found in popular textbooks, such as van der Waerden's Moderne Algebra, which start each chapter with a formal definition of a structure and then follow it with concrete examples.
Elementary algebra
The study of polynomial equations or algebraic equations has a long history. , the Babylonians were able to solve quadratic equations specified as word problems. This word problem stage is classified as rhetorical algebra and was the dominant approach up to the 16th century. Al-Khwarizmi originated the word "algebra" in 830 AD, but his work was entirely rhetorical algebra. Fully symbolic algebra did not appear until François Viète's 1591 New Algebra, and even this had some spelled out words that were given symbols in Descartes's 1637 La Géométrie. The formal study of solving symbolic equations led Leonhard Euler to accept what were then considered "nonsense" roots such as negative numbers and imaginary numbers, in the late 18th century. However, European mathematicians, for the most part, resisted these concepts until the middle of the 19th century.
George Peacock's 1830 Treatise of Algebra was the first attempt to place algebra on a strictly symbolic basis. He distinguished a new symbolical algebra, distinct from the old arithmetical algebra. Whereas in arithmetical algebra is restricted to , in symbolical algebra all rules of operations hold with no restrictions. Using this Peacock could show laws such as , by letting in . Peacock used what he termed the principle of the permanence of equivalent forms to justify his argument, but his reasoning suffered from the problem of induction. For example, holds for the nonnegative real numbers, but not for general complex numbers.
Early group theory
Several areas of mathematics led to the study of groups. Lagrange's 1770 study of the solutions of the quintic equation led to the Galois group of a polynomial. Gauss's 1801 study of Fermat's little theorem led to the ring of integers modulo n, the multiplicative group of integers modulo n, and the more general concepts of cyclic groups and abelian groups. Klein's 1872 Erlangen program studied geometry and led to symmetry groups such as the Euclidean group and the group of projective transformations. In 1874 Lie introduced the theory of Lie groups, aiming for "the Galois theory of differential equations". In 1876 Poincaré and Klein introduced the group of Möbius transformations, and its subgroups such as the modular group and Fuchsian group, based on work on automorphic functions in analysis.
The abstract concept of group emerged slowly over the middle of the nineteenth century. Galois in 1832 was the first to use the term "group", signifying a collection of permutations closed under composition. Arthur Cayley's 1854 paper On the theory of groups defined a group as a set with an associative composition operation and the identity 1, today called a monoid. In 1870 Kronecker defined an abstract binary operation that was closed, commutative, associative, and had the left cancellation property , similar to the modern laws for a finite abelian group. Weber's 1882 definition of a group was a closed binary operation that was associative and had left and right cancellation. Walther von Dyck in 1882 was the first to require inverse elements as part of the definition of a group.
Once this abstract group concept emerged, results were reformulated in this abstract setting. For example, Sylow's theorem was reproven by Frobenius in 1887 directly from the laws of a finite group, although Frobenius remarked that the theorem followed from Cauchy's theorem on permutation groups and the fact that every finite group is a subgroup of a permutation group. Otto Hölder was particularly prolific in this area, defining quotient groups in 1889, group automorphisms in 1893, as well as simple groups. He also completed the Jordan–Hölder theorem. Dedekind and Miller independently characterized Hamiltonian groups and introduced the notion of the commutator of two elements. Burnside, Frobenius, and Molien created the representation theory of finite groups at the end of the nineteenth century. J. A. de Séguier's 1905 monograph Elements of the Theory of Abstract Groups presented many of these results in an abstract, general form, relegating "concrete" groups to an appendix, although it was limited to finite groups. The first monograph on both finite and infinite abstract groups was O. K. Schmidt's 1916 Abstract Theory of Groups.
Early ring theory
Noncommutative ring theory began with extensions of the complex numbers to hypercomplex numbers, specifically William Rowan Hamilton's quaternions in 1843. Many other number systems followed shortly. In 1844, Hamilton presented biquaternions, Cayley introduced octonions, and Grassman introduced exterior algebras. James Cockle presented tessarines in 1848 and coquaternions in 1849. William Kingdon Clifford introduced split-biquaternions in 1873. In addition Cayley introduced group algebras over the real and complex numbers in 1854 and square matrices in two papers of 1855 and 1858.
Once there were sufficient examples, it remained to classify them. In an 1870 monograph, Benjamin Peirce classified the more than 150 hypercomplex number systems of dimension below 6, and gave an explicit definition of an associative algebra. He defined nilpotent and idempotent elements and proved that any algebra contains one or the other. He also defined the Peirce decomposition. Frobenius in 1878 and Charles Sanders Peirce in 1881 independently proved that the only finite-dimensional division algebras over were the real numbers, the complex numbers, and the quaternions. In the 1880s Killing and Cartan showed that semisimple Lie algebras could be decomposed into simple ones, and classified all simple Lie algebras. Inspired by this, in the 1890s Cartan, Frobenius, and Molien proved (independently) that a finite-dimensional associative algebra over or uniquely decomposes into the direct sums of a nilpotent algebra and a semisimple algebra that is the product of some number of simple algebras, square matrices over division algebras. Cartan was the first to define concepts such as direct sum and simple algebra, and these concepts proved quite influential. In 1907 Wedderburn extended Cartan's results to an arbitrary field, in what are now called the Wedderburn principal theorem and Artin–Wedderburn theorem.
For commutative rings, several areas together led to commutative ring theory. In two papers in 1828 and 1832, Gauss formulated the Gaussian integers and showed that they form a unique factorization domain (UFD) and proved the biquadratic reciprocity law. Jacobi and Eisenstein at around the same time proved a cubic reciprocity law for the Eisenstein integers. The study of Fermat's last theorem led to the algebraic integers. In 1847, Gabriel Lamé thought he had proven FLT, but his proof was faulty as he assumed all the cyclotomic fields were UFDs, yet as Kummer pointed out, was not a UFD. In 1846 and 1847 Kummer introduced ideal numbers and proved unique factorization into ideal primes for cyclotomic fields. Dedekind extended this in 1871 to show that every nonzero ideal in the domain of integers of an algebraic number field is a unique product of prime ideals, a precursor of the theory of Dedekind domains. Overall, Dedekind's work created the subject of algebraic number theory.
In the 1850s, Riemann introduced the fundamental concept of a Riemann surface. Riemann's methods relied on an assumption he called Dirichlet's principle, which in 1870 was questioned by Weierstrass. Much later, in 1900, Hilbert justified Riemann's approach by developing the direct method in the calculus of variations. In the 1860s and 1870s, Clebsch, Gordan, Brill, and especially M. Noether studied algebraic functions and curves. In particular, Noether studied what conditions were required for a polynomial to be an element of the ideal generated by two algebraic curves in the polynomial ring , although Noether did not use this modern language. In 1882 Dedekind and Weber, in analogy with Dedekind's earlier work on algebraic number theory, created a theory of algebraic function fields which allowed the first rigorous definition of a Riemann surface and a rigorous proof of the Riemann–Roch theorem. Kronecker in the 1880s, Hilbert in 1890, Lasker in 1905, and Macauley in 1913 further investigated the ideals of polynomial rings implicit in E. Noether's work. Lasker proved a special case of the Lasker-Noether theorem, namely that every ideal in a polynomial ring is a finite intersection of primary ideals. Macauley proved the uniqueness of this decomposition. Overall, this work led to the development of algebraic geometry.
In 1801 Gauss introduced binary quadratic forms over the integers and defined their equivalence. He further defined the discriminant of these forms, which is an invariant of a binary form. Between the 1860s and 1890s invariant theory developed and became a major field of algebra. Cayley, Sylvester, Gordan and others found the Jacobian and the Hessian for binary quartic forms and cubic forms. In 1868 Gordan proved that the graded algebra of invariants of a binary form over the complex numbers was finitely generated, i.e., has a basis. Hilbert wrote a thesis on invariants in 1885 and in 1890 showed that any form of any degree or number of variables has a basis. He extended this further in 1890 to Hilbert's basis theorem.
Once these theories had been developed, it was still several decades until an abstract ring concept emerged. The first axiomatic definition was given by Abraham Fraenkel in 1914. His definition was mainly the standard axioms: a set with two operations addition, which forms a group (not necessarily commutative), and multiplication, which is associative, distributes over addition, and has an identity element. In addition, he had two axioms on "regular elements" inspired by work on the p-adic numbers, which excluded now-common rings such as the ring of integers. These allowed Fraenkel to prove that addition was commutative. Fraenkel's work aimed to transfer Steinitz's 1910 definition of fields over to rings, but it was not connected with the existing work on concrete systems. Masazo Sono's 1917 definition was the first equivalent to the present one.
In 1920, Emmy Noether, in collaboration with W. Schmeidler, published a paper about the theory of ideals in which they defined left and right ideals in a ring. The following year she published a landmark paper called Idealtheorie in Ringbereichen (Ideal theory in rings), analyzing ascending chain conditions with regard to (mathematical) ideals. The publication gave rise to the term "Noetherian ring", and several other mathematical objects being called Noetherian. Noted algebraist Irving Kaplansky called this work "revolutionary"; results which seemed inextricably connected to properties of polynomial rings were shown to follow from a single axiom. Artin, inspired by Noether's work, came up with the descending chain condition. These definitions marked the birth of abstract ring theory.
Early field theory
In 1801 Gauss introduced the integers mod p, where p is a prime number. Galois extended this in 1830 to finite fields with elements. In 1871 Richard Dedekind introduced, for a set of real or complex numbers that is closed under the four arithmetic operations, the German word Körper, which means "body" or "corpus" (to suggest an organically closed entity). The English term "field" was introduced by Moore in 1893. In 1881 Leopold Kronecker defined what he called a domain of rationality, which is a field of rational fractions in modern terms. The first clear definition of an abstract field was due to Heinrich Martin Weber in 1893. It was missing the associative law for multiplication, but covered finite fields and the fields of algebraic number theory and algebraic geometry. In 1910 Steinitz synthesized the knowledge of abstract field theory accumulated so far. He axiomatically defined fields with the modern definition, classified them by their characteristic, and proved many theorems commonly seen today.
Other major areas
Solving of systems of linear equations, which led to linear algebra
Modern algebra
The end of the 19th and the beginning of the 20th century saw a shift in the methodology of mathematics. Abstract algebra emerged around the start of the 20th century, under the name modern algebra. Its study was part of the drive for more intellectual rigor in mathematics. Initially, the assumptions in classical algebra, on which the whole of mathematics (and major parts of the natural sciences) depend, took the form of axiomatic systems. No longer satisfied with establishing properties of concrete objects, mathematicians started to turn their attention to general theory. Formal definitions of certain algebraic structures began to emerge in the 19th century. For example, results about various groups of permutations came to be seen as instances of general theorems that concern a general notion of an abstract group. Questions of structure and classification of various mathematical objects came to forefront.
These processes were occurring throughout all of mathematics, but became especially pronounced in algebra. Formal definition through primitive operations and axioms were proposed for many basic algebraic structures, such as groups, rings, and fields. Hence such things as group theory and ring theory took their places in pure mathematics. The algebraic investigations of general fields by Ernst Steinitz and of commutative and then general rings by David Hilbert, Emil Artin and Emmy Noether, building on the work of Ernst Kummer, Leopold Kronecker and Richard Dedekind, who had considered ideals in commutative rings, and of Georg Frobenius and Issai Schur, concerning representation theory of groups, came to define abstract algebra. These developments of the last quarter of the 19th century and the first quarter of 20th century were systematically exposed in Bartel van der Waerden's Moderne Algebra, the two-volume monograph published in 1930–1931 that reoriented the idea of algebra from the theory of equations to the theory of algebraic structures.
Basic concepts
By abstracting away various amounts of detail, mathematicians have defined various algebraic structures that are used in many areas of mathematics. For instance, almost all systems studied are sets, to which the theorems of set theory apply. Those sets that have a certain binary operation defined on them form magmas, to which the concepts concerning magmas, as well those concerning sets, apply. We can add additional constraints on the algebraic structure, such as associativity (to form semigroups); identity, and inverses (to form groups); and other more complex structures. With additional structure, more theorems could be proved, but the generality is reduced. The "hierarchy" of algebraic objects (in terms of generality) creates a hierarchy of the corresponding theories: for instance, the theorems of group theory may be used when studying rings (algebraic objects that have two binary operations with certain axioms) since a ring is a group over one of its operations. In general there is a balance between the amount of generality and the richness of the theory: more general structures have usually fewer nontrivial theorems and fewer applications.
Examples of algebraic structures with a single binary operation are:
Magma
Quasigroup
Monoid
Semigroup
Group
Examples involving several operations include:
Ring
Field
Module
Vector space
Algebra over a field
Associative algebra
Lie algebra
Lattice
Boolean algebra
Branches of abstract algebra
Group theory
A group is a set together with a "group product", a binary operation . The group satisfies the following defining axioms (c.f. ):Identity: there exists an element such that, for each element in , it holds that .Inverse: for each element of , there exists an element so that .Associativity''': for each triplet of elements in , it holds that .
Ring theory
A ring is a set with two binary operations, addition: and multiplication: satisfying the following axioms.
is a commutative group under addition.
is a monoid under multiplication.
Multiplication is distributive with respect to addition.
Applications
Because of its generality, abstract algebra is used in many fields of mathematics and science. For instance, algebraic topology uses algebraic objects to study topologies. The Poincaré conjecture, proved in 2003, asserts that the fundamental group of a manifold, which encodes information about connectedness, can be used to determine whether a manifold is a sphere or not. Algebraic number theory studies various number rings that generalize the set of integers. Using tools of algebraic number theory, Andrew Wiles proved Fermat's Last Theorem.
In physics, groups are used to represent symmetry operations, and the usage of group theory could simplify differential equations. In gauge theory, the requirement of local symmetry can be used to deduce the equations describing a system. The groups that describe those symmetries are Lie groups, and the study of Lie groups and Lie algebras reveals much about the physical system; for instance, the number of force carriers in a theory is equal to the dimension of the Lie algebra, and these bosons interact with the force they mediate if the Lie algebra is nonabelian.
| Mathematics | Algebra | null |
19620261 | https://en.wikipedia.org/wiki/Ur%20%28continent%29 | Ur (continent) | Ur is a hypothetical supercontinent that formed in the Archean eon around 3.1 billion years ago (Ga). In a reconstruction by Rogers, Ur is half a billion years older than Arctica and, in the early period of its existence, probably the only continent on Earth, making it a supercontinent despite probably being smaller than present-day Australia. In more recent works geologists often refer to both Ur and other proposed Archaean continental assemblages as supercratons. Ur can, nevertheless, be half a billion years younger than Vaalbara, but the concepts of these two early cratonic assemblages are incompatible.
Incompatible reconstructions
About 1.3–1.1 Ga, Ur joined the continents Nena and Atlantica to form the supercontinent Rodinia. In the reconstruction of , Ur remained the nucleus of eastern Gondwana until the break-up of Gondwana. In other reconstructions, however, India and East Antarctica did not collide until Rodinia formed 1.1 Ga. Furthermore, in the early Archaean Earth's mantle was 200 °C hotter than today, and many characteristics of modern tectonics, such as ophiolites, blueschists, lawsonite-bearing eclogites, and ultra-high-pressure rocks, did not exist or were rare. This makes most proposed Archaean supercontinents controversial, including Rogers's 3 Ga supercontinent.
Reconstructions of Vaalbara place two cratons—Kaapvaal in southern Africa and Pilbara in western Australia—next to each other based on stratigraphic similarities. In Roger's configuration of Ur, these cratons are placed far apart in their Gondwana configuration. This configuration is contradicted by widespread Precambrian collisional events between Australia and Africa.
Another possible supercraton, Zimgarn, proposed by and named after the Zimbabwe and Yilgarn cratons, is distinct from both Vaalbara and Ur. Vaalbara and Zimgarn, according to this proposal, both disintegrated about 2.1–2.0 Ga to reassemble as the Kalahari and West Australian cratons 2.5–1.5 Ga. Smirnov et al. based this reconstruction on: (1) Zimgarn was still undergoing cratonisation when an extensive carbonate platform developed over Vaalbara; (2) the magmatic signatures are different for the two supercratons during the period 2.6–2.0 Ga; and (3) paleomagnetic latitudes for 2.7 Ga are slightly different.
Original concept and later interpretations
Important geological similarities link now remote Archaean cratons in India (Singhbhum and Dharwar), western Australia (Kilbaran and Pilbara), and southern Africa (Kaapvaal and Zimbabwe) which indicate that these protolithic shields were close together in the mid-Archaean. The name "Ur", from the German prefix ur- meaning "original", was introduced by , since it is the first continent in his tectonic reconstructions. Other Archaean continental assemblages are considerably younger: Arctica (Baltica, Laurentia, Ur, and Siberia) consolidated around 2.6 Ga, Atlantica (West Africa and eastern South America) consolidated around 2.1 Ga. In some reconstructions the shields of Ur stayed near each other until the Mesozoic break-up of Gondwana.
The cratons that had become stable around 3 Ga were all in the same region within Pangaea, which is the main argument for them having formed a single continent 3 Ga. The Kaapvaal Craton became stable around 3.1 Ga. The Pilbara Craton is not well defined but formed around 3 Ga. Three cratons in East Antarctica are of similar age but not well known. These cratons share similar geological histories and are therefore assumed to have formed Vaalbara.
Three small areas in the Indian Ocean coast of Antarctica are also about 3 Ga old: western Queen Maud Land, the Napier complex, and the Vestfold Hills. Within Gondwana, these areas were in a belt of Grenville-age deformation, and because there is no evidence of ocean closure in this belt (except in Africa), the 1 Ga orogen can be assumed to be intra-continental. Consequently, the southern margin of Ur is now below the Antarctic ice cover.
Two cratons in India of equal age, Western Dharwar and Singhbhum, were also part of Ur. Two other Indian cratons, Eastern Dharwar and Bhandara, also formed around 3 Ga but underwent extensive magmatism around 2.5 Ga not seen elsewhere, and their relation to Ur is unclear. Ur, nevertheless, became larger around 2.5 Ga, and this so-called "Expanded Ur" incorporated the Zimbabwe and Yilgarn cratons.
The largest preserved parts of Ur are in India: Aravalli, Dharwar, Bundelkhand, and Singhbhum. The Central Indian Tectonic Zone is the modern suture between the Bundelkhand-Aravalli block and the other Archaean blocks. 2.8–2.6 Ga metamorphism in Dharwar and Bundelkhand indicate that the stabilisation of Ur probably continued until the end of that period.
| Physical sciences | Paleogeography | Earth science |
259906 | https://en.wikipedia.org/wiki/Theorema%20Egregium | Theorema Egregium | Gauss's Theorema Egregium (Latin for "Remarkable Theorem") is a major result of differential geometry, proved by Carl Friedrich Gauss in 1827, that concerns the curvature of surfaces. The theorem says that Gaussian curvature can be determined entirely by measuring angles, distances and their rates on a surface, without reference to the particular manner in which the surface is embedded in the ambient 3-dimensional Euclidean space. In other words, the Gaussian curvature of a surface does not change if one bends the surface without stretching it. Thus the Gaussian curvature is an intrinsic invariant of a surface.
Gauss presented the theorem in this manner (translated from Latin):
Thus the formula of the preceding article leads itself to the remarkable Theorem. If a curved surface is developed upon any other surface whatever, the measure of curvature in each point remains unchanged.
The theorem is "remarkable" because the definition of Gaussian curvature makes ample reference to the specific way the surface is embedded in 3-dimensional space, and it is quite surprising that the result does not depend on its embedding.
In modern mathematical terminology, the theorem may be stated as follows:
Elementary applications
A sphere of radius R has constant Gaussian curvature which is equal to 1/R2. At the same time, a plane has zero Gaussian curvature. As a corollary of Theorema Egregium, a piece of paper cannot be bent onto a sphere without crumpling. Conversely, the surface of a sphere cannot be unfolded onto a flat plane without distorting the distances. If one were to step on an empty egg shell, its edges have to split in expansion before being flattened. Mathematically, a sphere and a plane are not isometric, even locally. This fact is significant for cartography: it implies that no planar (flat) map of Earth can be perfect, even for a portion of the Earth's surface. Thus every cartographic projection necessarily distorts at least some distances.
The catenoid and the helicoid are two very different-looking surfaces. Nevertheless, each of them can be continuously bent into the other: they are locally isometric. It follows from Theorema Egregium that under this bending the Gaussian curvature at any two corresponding points of the catenoid and helicoid is always the same. Thus isometry is simply bending and twisting of a surface without internal crumpling or tearing, in other words without extra tension, compression, or shear.
An application of the theorem is seen when a flat object is somewhat folded or bent along a line, creating rigidity in the perpendicular direction. This is of practical use in construction, as well as in a common pizza-eating strategy: A flat slice of pizza can be seen as a surface with constant Gaussian curvature 0. Gently bending a slice must then roughly maintain this curvature (assuming the bend is roughly a local isometry). If one bends a slice horizontally along a radius, non-zero principal curvatures are created along the bend, dictating that the other principal curvature at these points must be zero. This creates rigidity in the direction perpendicular to the fold, an attribute desirable for eating pizza, as it holds its shape long enough to be consumed without a mess. This same principle is used for strengthening in corrugated materials, most familiarly with corrugated fiberboard and corrugated galvanised iron, and in some forms of potato chips as well.
| Mathematics | Other | null |
259962 | https://en.wikipedia.org/wiki/Avoirdupois | Avoirdupois | Avoirdupois (; abbreviated avdp.) is a measurement system of weights that uses pounds and ounces as units.<ref
name=sizes /> It was first commonly used in the 13th century AD and was updated in 1959.
In 1959, by international agreement, the definitions of the pound and ounce became standardized in countries which use the pound as a unit of mass. The International Avoirdupois Pound was then created. It is the everyday system of weights used in the United States. It is still used, in varying degrees, in everyday life in the United Kingdom, Canada, Australia, and some other former British colonies, despite their official adoption of the metric system.
The avoirdupois weight system's general attributes were originally developed for the international wool trade in the Late Middle Ages, when trade was in recovery. It was historically based on a physical standardized pound or "prototype weight" that could be divided into 16 ounces. There were a number of competing measures of mass, and the fact that the avoirdupois pound had three even numbers as divisors (half and half and half again) may have been a cause of much of its popularity, so that the system won out over systems with 12 or 10 or 15 subdivisions. The use of this unofficial system gradually stabilized and evolved, with only slight changes in the reference standard or in the prototype's actual mass.Over time, the desire not to use too many different systems of measurement allowed the establishment of "value relationships", with other commodities metered and sold by weight measurements such as bulk goods (grains, ores, flax) and smelted metals; so the avoirdupois system gradually became an accepted standard through much of Europe.
In England, Henry VII authorized its use as a standard, and Queen Elizabeth I acted three times to enforce a common standard, thus establishing what became the Imperial system of weights and measures. Late in the 19th century various governments acted to redefine their base standards on a scientific basis and establish ratios between local avoirdupois measurements and international SI metric system standards. The legal actions of these various governments were independently conceived, and so did not always pick the same ratios to metric units for each avoirdupois unit. The result of this was, after these standardisations, measurements of the same name often had marginally different recognised values in different regions (although the pound generally remained very similar). In the modern day, this is evident in the small difference between United States customary and British Imperial pounds.
An alternative system of mass, the troy system, is generally used for precious materials. The modern definition of the avoirdupois pound (1 lb) is exactly kilograms.
Etymology
is from Anglo-Norman French (later ), literally "goods of weight" (Old French , as verb meaning "to have" and as noun meaning "property, goods", comes from the Latin , "to have, to hold, to possess something"; = "from"/"of", cf. Latin; = "weight", from Latin .) This term originally referred to a class of merchandise: , "goods of weight", things that were sold in bulk and were weighed on large steelyards or balances.
Only later did the term become identified with a particular system of units used to weigh such merchandise. Inconsistent orthography throughout history has left many variants of the term, such as and . (The Norman became the Parisian . In the 17th century was replaced with .)
The current spelling of the last word is in the current standard French orthography, but the spelling avoirdupois remained as is in the anglosphere.
History
The rise in use of the measurement system corresponds to the regrowth of trade during the High Middle Ages after the early crusades, when Europe experienced a growth in towns, turned from the chaos of warlordism to long-distance trade, and began annual fairs, tournaments and commerce, by land and sea. There are two major hypotheses regarding the origins of the avoirdupois system. The older hypothesis is that it originated in France. A newer hypothesis is that it is based on the weight system of Florence.
The avoirdupois weight system is thought to have come into use in England around 1300. It was originally used for weighing wool. In the early 14th century several other specialized weight systems were used, including the weight system of the Hanseatic League with a 16-ounce pound of grains and an 8-ounce mark. However, the main weight system, used for coinage and for everyday use, was based on the 12-ounce tower pound of grains. From the 14th century until the late 16th century, the system's basis and the prototype for today's international pound, the avoirdupois pound, was also known as the wool pound or the avoirdupois wool pound.
The earliest known version of the avoirdupois weight system had the following units: a pound of grains, a stone of 14 pounds, a woolsack of 26 stone, an ounce of pound, and finally, the ounce was divided into 16 "parts".
The earliest known occurrence of the word "avoirdupois" (or some variant thereof) in England is from a document entitled Tractatus de Ponderibus et Mensuris ("Treatise on Weights and Measures"). This document is listed in early statute books under the heading 31 Edward I dated 2 February 1303. More recent statute books list it among statutes of uncertain date. Scholars nowadays believe that it was probably written between 1266 and 1303. Initially a royal memorandum, it eventually took on the force of law and was recognized as a statute by King Henry VIII and Queen Elizabeth I.
It was repealed by the Weights and Measures Act 1824 (5 Geo. 4. c. 74). In the Tractatus, the word "avoirdupois" refers not to a weight system, but to a class of goods, specifically heavy goods sold by weight, as opposed to goods sold by volume, count, or some other method. Since it is written in Anglo-Norman French, this document is not the first occurrence of the word in the English language.
Toward a uniformity of measures
Three major developments occurred during the reign of Edward III (r. 1327–1377). First, a statute cited as 14 Edw 3 Stat. 1 c. 12 (1340) "Bushels and Weights shall be made and sent into every County."
The second major development is the statute 25 Edw 3 Stat. 5 c. 9 (1350) "The Auncel Weight shall be put out, and Weighing shall be by equal Balance."
The third development is a set of 14th-century bronze weights at the Westgate Museum in Winchester, England. The weights are in denominations of 7 pounds (corresponding to a unit known as the clip or wool-clip), 14 pounds (stone), 56 pounds (4 stone) and 91 pounds ( sack or woolsack). The 91-pound weight is thought to have been commissioned by Edward III in conjunction with the statute of 1350, while the other weights are thought to have been commissioned in conjunction with the statutes of 1340. The 56-pound weight was used as a reference standard as late as 1588.
A statute of Henry VIII (24 Hen 8 c. 3) made avoirdupois weights mandatory.
In 1588 Queen Elizabeth increased the weight of the avoirdupois pound to grains and added the troy grain to the avoirdupois weight system. Prior to 1588, the "part" () was the smallest unit in the avoirdupois weight system. In the 18th century, the "part" was renamed "drachm".
Original forms
These are the units in their original Anglo-Norman French forms:
Post-Elizabethan units
In the United Kingdom, 14 avoirdupois pounds equal one stone. The quarter, hundredweight, and ton equal respectively, 28 lb, 112 lb, and 2,240 lb in order for masses to be easily converted between them and stone. The following are the units in the British or imperial version of the avoirdupois system:
Note: The plural form of the unit stone is either stone or stones, but stone is most frequently used.
American customary system
The thirteen British colonies in North America used the avoirdupois system, but continued to use the British system as it was, without the evolution that was occurring in Britain in the use of the stone unit. In 1824 there was landmark new weights and measures legislation in the United Kingdom that the United States did not adopt.
In the United States, quarters, hundredweights, and tons remain defined as 25, 100, and respectively. The quarter is now virtually unused, as is the hundredweight outside of agriculture and commodities. If disambiguation is required, then they are referred to as the smaller "short" units in the United States, as opposed to the larger British "long" units. Grains are used worldwide for measuring gunpowder and smokeless powder charges. Historically, the dram ( grains; not to be mixed up with the apothecaries' dram of 60 grains) has also been used worldwide for measuring gunpowder charges, particularly for shotguns and large black-powder rifles.
| Physical sciences | Measurement systems | Basics and measurement |
778700 | https://en.wikipedia.org/wiki/Laws%20of%20thermodynamics | Laws of thermodynamics | The laws of thermodynamics are a set of scientific laws which define a group of physical quantities, such as temperature, energy, and entropy, that characterize thermodynamic systems in thermodynamic equilibrium. The laws also use various parameters for thermodynamic processes, such as thermodynamic work and heat, and establish relationships between them. They state empirical facts that form a basis of precluding the possibility of certain phenomena, such as perpetual motion. In addition to their use in thermodynamics, they are important fundamental laws of physics in general and are applicable in other natural sciences.
Traditionally, thermodynamics has recognized three fundamental laws, simply named by an ordinal identification, the first law, the second law, and the third law. A more fundamental statement was later labelled as the zeroth law after the first three laws had been established.
The zeroth law of thermodynamics defines thermal equilibrium and forms a basis for the definition of temperature: if two systems are each in thermal equilibrium with a third system, then they are in thermal equilibrium with each other.
The first law of thermodynamics states that, when energy passes into or out of a system (as work, heat, or matter), the system's internal energy changes in accordance with the law of conservation of energy.
The second law of thermodynamics states that in a natural thermodynamic process, the sum of the entropies of the interacting thermodynamic systems never decreases. A common corollary of the statement is that heat does not spontaneously pass from a colder body to a warmer body.
The third law of thermodynamics states that a system's entropy approaches a constant value as the temperature approaches absolute zero. With the exception of non-crystalline solids (glasses), the entropy of a system at absolute zero is typically close to zero.
The first and second laws prohibit two kinds of perpetual motion machines, respectively: the perpetual motion machine of the first kind which produces work with no energy input, and the perpetual motion machine of the second kind which spontaneously converts thermal energy into mechanical work.
History
The history of thermodynamics is fundamentally interwoven with the history of physics and the history of chemistry, and ultimately dates back to theories of heat in antiquity. The laws of thermodynamics are the result of progress made in this field over the nineteenth and early twentieth centuries. The first established thermodynamic principle, which eventually became the second law of thermodynamics, was formulated by Sadi Carnot in 1824 in his book Reflections on the Motive Power of Fire. By 1860, as formalized in the works of scientists such as Rudolf Clausius and William Thomson, what are now known as the first and second laws were established. Later, Nernst's theorem (or Nernst's postulate), which is now known as the third law, was formulated by Walther Nernst over the period 1906–1912. While the numbering of the laws is universal today, various textbooks throughout the 20th century have numbered the laws differently. In some fields, the second law was considered to deal with the efficiency of heat engines only, whereas what was called the third law dealt with entropy increases. Gradually, this resolved itself and a zeroth law was later added to allow for a self-consistent definition of temperature. Additional laws have been suggested, but have not achieved the generality of the four accepted laws, and are generally not discussed in standard textbooks.
Zeroth law
The zeroth law of thermodynamics provides for the foundation of temperature as an empirical parameter in thermodynamic systems and establishes the transitive relation between the temperatures of multiple bodies in thermal equilibrium. The law may be stated in the following form:
Though this version of the law is one of the most commonly stated versions, it is only one of a diversity of statements that are labeled as "the zeroth law". Some statements go further, so as to supply the important physical fact that temperature is one-dimensional and that one can conceptually arrange bodies in a real number sequence from colder to hotter.
These concepts of temperature and of thermal equilibrium are fundamental to thermodynamics and were clearly stated in the nineteenth century. The name 'zeroth law' was invented by Ralph H. Fowler in the 1930s, long after the first, second, and third laws were widely recognized. The law allows the definition of temperature in a non-circular way without reference to entropy, its conjugate variable. Such a temperature definition is said to be 'empirical'.
First law
The first law of thermodynamics is a version of the law of conservation of energy, adapted for thermodynamic processes. In general, the conservation law states that the total energy of an isolated system is constant; energy can be transformed from one form to another, but can be neither created nor destroyed.
For processes that include the transfer of matter, a further statement is needed.
The First Law encompasses several principles:
Conservation of energy, which says that energy can be neither created nor destroyed, but can only change form. A particular consequence of this is that the total energy of an isolated system does not change.
The concept of internal energy and its relationship to temperature. If a system has a definite temperature, then its total energy has three distinguishable components, termed kinetic energy (energy due to the motion of the system as a whole), potential energy (energy resulting from an externally imposed force field), and internal energy. The establishment of the concept of internal energy distinguishes the first law of thermodynamics from the more general law of conservation of energy.
Work is a process of transferring energy to or from a system in ways that can be described by macroscopic mechanical forces acting between the system and its surroundings. The work done by the system can come from its overall kinetic energy, from its overall potential energy, or from its internal energy. For example, when a machine (not a part of the system) lifts a system upwards, some energy is transferred from the machine to the system. The system's energy increases as work is done on the system and in this particular case, the energy increase of the system is manifested as an increase in the system's gravitational potential energy. Work added to the system increases the potential energy of the system.
When matter is transferred into a system, the internal energy and potential energy associated with it are transferred into the new combined system. where denotes the internal energy per unit mass of the transferred matter, as measured while in the surroundings; and denotes the amount of transferred mass.
The flow of heat is a form of energy transfer. Heat transfer is the natural process of moving energy to or from a system, other than by work or the transfer of matter. In a diathermal system, the internal energy can only be changed by the transfer of energy as heat:
Combining these principles leads to one traditional statement of the first law of thermodynamics: it is not possible to construct a machine which will perpetually output work without an equal amount of energy input to that machine. Or more briefly, a perpetual motion machine of the first kind is impossible.
Second law
The second law of thermodynamics indicates the irreversibility of natural processes, and in many cases, the tendency of natural processes to lead towards spatial homogeneity of matter and energy, especially of temperature. It can be formulated in a variety of interesting and important ways. One of the simplest is the Clausius statement, that heat does not spontaneously pass from a colder to a hotter body.
It implies the existence of a quantity called the entropy of a thermodynamic system. In terms of this quantity it implies that
The second law is applicable to a wide variety of processes, both reversible and irreversible. According to the second law, in a reversible heat transfer, an element of heat transferred, , is the product of the temperature (), both of the system and of the sources or destination of the heat, with the increment () of the system's conjugate variable, its entropy ():
While reversible processes are a useful and convenient theoretical limiting case, all natural processes are irreversible. A prime example of this irreversibility is the transfer of heat by conduction or radiation. It was known long before the discovery of the notion of entropy that when two bodies, initially of different temperatures, come into direct thermal connection, then heat immediately and spontaneously flows from the hotter body to the colder one.
Entropy may also be viewed as a physical measure concerning the microscopic details of the motion and configuration of a system, when only the macroscopic states are known. Such details are often referred to as disorder on a microscopic or molecular scale, and less often as dispersal of energy. For two given macroscopically specified states of a system, there is a mathematically defined quantity called the 'difference of information entropy between them'. This defines how much additional microscopic physical information is needed to specify one of the macroscopically specified states, given the macroscopic specification of the other – often a conveniently chosen reference state which may be presupposed to exist rather than explicitly stated. A final condition of a natural process always contains microscopically specifiable effects which are not fully and exactly predictable from the macroscopic specification of the initial condition of the process. This is why entropy increases in natural processes – the increase tells how much extra microscopic information is needed to distinguish the initial macroscopically specified state from the final macroscopically specified state. Equivalently, in a thermodynamic process, energy spreads.
Third law
The third law of thermodynamics can be stated as:
At absolute zero temperature, the system is in the state with the minimum thermal energy, the ground state. The constant value (not necessarily zero) of entropy at this point is called the residual entropy of the system. With the exception of non-crystalline solids (e.g. glass) the residual entropy of a system is typically close to zero. However, it reaches zero only when the system has a unique ground state (i.e., the state with the minimum thermal energy has only one configuration, or microstate). Microstates are used here to describe the probability of a system being in a specific state, as each microstate is assumed to have the same probability of occurring, so macroscopic states with fewer microstates are less probable. In general, entropy is related to the number of possible microstates according to the Boltzmann principle
where S is the entropy of the system, kB is the Boltzmann constant, and Ω the number of microstates. At absolute zero there is only 1 microstate possible (Ω = 1 as all the atoms are identical for a pure substance, and as a result all orders are identical as there is only one combination) and .
Onsager relations
The Onsager reciprocal relations have been considered the fourth law of thermodynamics. They describe the relation between thermodynamic flows and forces in non-equilibrium thermodynamics, under the assumption that thermodynamic variables can be defined locally in a condition of local equilibrium. These relations are derived from statistical mechanics under the principle of microscopic reversibility (in the absence of external magnetic fields). Given a set of extensive parameters (energy, mass, entropy, number of particles and so on) and thermodynamic forces (related to their related intrinsic parameters, such as temperature and pressure), the Onsager theorem states that
where index every parameter and its related force, and
are called the thermodynamic flows.
| Physical sciences | Thermodynamics | null |
779725 | https://en.wikipedia.org/wiki/Red%20deer | Red deer | The red deer (Cervus elaphus) is one of the largest deer species. A male red deer is called a stag or hart, and a female is called a doe or hind. The red deer inhabits most of Europe, the Caucasus Mountains region, Anatolia, Iran, and parts of western Asia. It also inhabits the Atlas Mountains of Northern Africa; being the only living species of deer to inhabit Africa. Red deer have been introduced to other areas, including Australia, New Zealand, the United States, Canada, Peru, Uruguay, Chile and Argentina. In many parts of the world, the meat (venison) from red deer is used as a food source.
The red deer is a ruminant, characterized by a four-chambered stomach. Genetic evidence indicates that the red deer, as traditionally defined, is a species group, rather than a single species, though exactly how many species the group includes remains disputed. The ancestor of the red deer probably originated in central Asia.
Although at one time red deer were rare in parts of Europe, they were never close to extinction. Reintroduction and conservation efforts, such as in the United Kingdom and Portugal, have resulted in an increase of red deer populations, while other areas, such as North Africa, have continued to show a population decline.
Description
The red deer is the fourth-largest extant deer species, behind the moose, elk, and sambar deer. It is a ruminant, eating its food in two stages and having an even number of toes on each hoof, like camels, goats, and cattle. European red deer have a relatively long tail compared with their Asian and North American relatives. Subtle differences in appearance are noted between the various subspecies of red deer, primarily in size and antlers, with the smallest being the Corsican red deer found on the islands of Corsica and Sardinia and the largest being the Caspian red deer (or maral) of Asia Minor and the Caucasus Region to the west of the Caspian Sea.
The deer of central and western Europe vary greatly in size, with some of the largest deer found in the Carpathian Mountains in Central Europe. Western European red deer, historically, grew to large size given ample food supply (including people's crops), and descendants of introduced populations living in New Zealand and Argentina have grown quite large in both body and antler size. Large red deer stags, like the Caspian red deer or those of the Carpathian Mountains, may rival North American elk in size. Female red deer are much smaller than their male counterparts.
Size
The male (stag) red deer is typically long from the nose to the base of the tail and typically weighs ; the female (hind) is long and often weighs . The tail adds another and shoulder height is about . In Scotland, stags average in head-and-body length and high at the shoulder and females average long and tall. Based on body mass, they are likely the fourth largest extant deer species on average, behind the moose, the elk and the sambar deer.
Size varies in different subspecies with the largest, the huge but small-antlered deer of the Carpathian Mountains (C. e. elaphus), weighing up to . At the other end of the scale, the Corsican red deer (C. e. corsicanus) weighs about , although red deer in poor habitats can weigh as little as .
Neck mane
The males of many subspecies also grow a short neck mane during the autumn. The male deer of the British Isles and Norway tend to have the thickest and most noticeable manes. Male Caspian red deer (C. e. maral) and Spanish red deer (C. e. hispanicus) do not carry neck manes. Male deer of all subspecies, however, tend to have stronger and thicker neck muscles than female deer, which may give them an appearance of having neck manes. Red deer hinds (females) do not have neck manes.
Antlers
Only the stags have antlers, which start growing in the spring and are shed each year, usually at the end of winter. Antlers typically measure in total length and weigh , although large ones can grow to and weigh . Antlers, which are made of bone, can grow at a rate of a day. While an antler is growing, it is covered with highly vascular skin called velvet, which supplies oxygen and nutrients to the growing bone.
The antlers are testosterone-driven and as the stag's testosterone levels drop in the autumn, the velvet is shed and the antlers stop growing. With the approach of autumn, the antlers begin to calcify and the stags' testosterone production builds for the approaching rut (mating season).
European red deer antlers are distinctive in being rather straight and rugose, with the fourth and fifth tines forming a "crown" or "cup" in larger males. Any tines in excess of the fourth and fifth tines grow radially from the cup, which are generally absent in the antlers of smaller red deer, such as Corsican red deer. Western European red deer antlers feature "bez" (second) tines that are either absent or smaller than the brow tines. However, bez tines occur frequently in Norwegian red deer. Antlers of Caspian red deer carry large bez tines and form less-developed cups than western European red deer, their antlers are thus more like the "throw back" top tines of the North American elk (C. canadensis), known as maraloid characteristics. A stag can (exceptionally) have antlers with no tines, and is then known as a switch. Similarly, a stag that does not grow antlers is a hummel.
Coat
European red deer tend to be reddish-brown in their summer coats, and some individuals may have a few spots on the backs of their summer coats. During the autumn, all red deer subspecies grow thicker coats of hair, which helps to insulate them during the winter. Autumn is also when some of the stags grow their neck manes. The autumn/winter coats of most subspecies are most distinct. The Caspian red deer's winter coat is greyer and has a larger and more distinguished light rump-patch (like wapiti and some central Asian red deer) compared with the Western European red deer, which has more of a greyish-brown coat with a darker yellowish rump patch in the winter.
By the time summer begins, the heavy winter coat has been shed; the animals are known to rub against trees and other objects to help remove hair from their bodies. Red deer have different colouration based on the seasons and types of habitats, with grey or lighter colouration prevalent in the winter and more reddish and darker coat colouration in the summer.
Distribution
Europe and North Africa
The European red deer is found in southwestern Asia (Asia Minor and Caucasus regions), North Africa, and Europe. The red deer is the largest nondomesticated land mammal still existing in Ireland. The Barbary stag (which resembles the western European red deer) is the only living member of the deer family native to Africa, with the population centred in the northwestern region of the continent in the Atlas Mountains. As of the mid-1990s, Morocco, Tunisia, and Algeria were the only African countries known to have red deer.
In the Netherlands, a large herd (about 3000 animals counted in late 2012) lives in the Oostvaardersplassen, a nature reserve. Ireland has its own unique subspecies. In France, the population is thriving, having multiplied five-fold in the last half-century, increasing from 30,000 in 1970 to around 160,000 in 2014. The deer has particularly expanded its footprint into forests at higher altitudes than before.
In the UK, indigenous populations occur in Scotland, the Lake District, and the south west of England (principally on Exmoor). Not all of these are of entirely pure bloodlines, as some of these populations have been supplemented with deliberate releases of deer from parks, such as Warnham or Woburn Abbey, in an attempt to increase antler sizes and body weights. The University of Edinburgh found that, in Scotland, extensive hybridisation with the closely related sika deer has occurred.
Several other populations have originated either with "carted" deer kept for stag hunts being left out at the end of the hunt, escapes from deer farms, or deliberate releases. Carted deer were kept by stag hunts with no wild red deer in the locality and were normally recaptured after the hunt and used again; although the hunts are called "stag hunts", the Norwich Staghounds only hunted hinds (female red deer), and in 1950, at least eight hinds (some of which may have been pregnant) were known to be at large near Kimberley and West Harling; they formed the basis of a new population based in Thetford Forest in Norfolk. Further substantial red deer herds originated from escapes or deliberate releases in the New Forest, the Peak District, Suffolk, Lancashire, Brecon Beacons, and North Yorkshire, as well as many other smaller populations scattered throughout England and Wales, and they are all generally increasing in numbers and range. A census of deer populations in 2007 and again in 2011 coordinated by the British Deer Society records the red deer as having continued to expand their range in England and Wales since 2000, with expansion most notable in the Midlands and East Anglia.
Iran
Caspian red deer are found in the Hyrcanian Forests.
New Zealand
In New Zealand, red deer were introduced by acclimatisation societies along with other deer and game species. The first red deer to reach New Zealand were a pair sent by Lord Petre in 1851 from his herd at Thorndon Park, Essex, to the South Island, but the hind was shot before they had a chance to breed. Lord Petre sent another stag and two hinds in 1861, and these were liberated near Nelson, from where they quickly spread. The first deer to reach the North Island were a gift to Sir Frederick Weld from Windsor Great Park and were released near Wellington; these were followed by further releases up to 1914. Between 1851 and 1926, 220 separate liberations of red deer involved over 800 deer. In 1927, the State Forest Service introduced a bounty for red deer shot on their land, and in 1931, government control operations were commenced. Between 1931 and March 1975, 1,124,297 deer were killed on official operations.
The introduced red deer have adapted well and are widely hunted on both islands; many of the 220 introductions used deer originating from Scotland (Invermark) or one of the major deer parks in England, principally Warnham, Woburn Abbey or Windsor Great Park. Some hybridisation happened with the closely related American elk (Cervus canadensis nelsoni) introduced in Fiordland in 1921. Along with the other introduced deer species, they are, however, officially regarded as a noxious pest and are still heavily culled using professional hunters working with helicopters, or even poisoned.
Australia
The first red deer to reach Australia were probably the six that Prince Albert sent in 1860 from Windsor Great Park to Thomas Chirnside, who was starting a herd at Werribee Park, south west of Melbourne in Victoria. Further introductions were made in New South Wales, Queensland, South Australia, and Western Australia. Today, red deer in Australia range from Queensland south through New South Wales into Victoria and across to South Australia, with the numbers increasing. The Queensland, Victorian and most New South Wales strains can still be traced to the early releases, but South Australia's population, along with all others, is now largely recent farm escapees. This is having adverse effects on the integrity of wild herds, as now more and larger herds are being grown due to the superior genetics that have been attained by selective breeding.
Wild red deer are a feral pest species in Australia, do considerable harm to the natural environment, and are a significant road traffic hazard.
Argentina and Chile
In Argentina and Chile, the red deer has had a potentially adverse impact on native animal species, such as the South Andean deer or huemul; the International Union for Conservation of Nature and Natural Resources has labelled the animal as one of the world's 100 worst invaders.
Migration
Red deer in Europe generally spend their winters at lower altitudes in more wooded terrain. During the summer, they migrate to higher elevations where food supplies are greater and better for the calving season.
Taxonomy and evolution
Until recently, biologists considered the red deer and elk or wapiti (C. canadensis) the same species, forming a continuous distribution throughout temperate Eurasia and North America. This belief was based largely on the fully fertile hybrids that can be produced under captive conditions.
Genetic evidence clearly shows the wapiti and red deer form two separate species.
Another member of the red deer group which may represent a separate species is C. corsicanus. If so, C. corsicanus includes the subspecies C. e. barbarus (perhaps a synonym of C. e. corsicanus), and is restricted to Maghreb in North Africa, Corsica, and Sardinia.
A 2014 mitochondrial DNA study showed the internal phylogeny of Cervus to be as follows:
Cervus elaphus appeared in Europe by the beginning of the Middle Pleistocene around 800,000 years ago. These earliest forms belonged to the palaeosubspecies Cervus elaphus acoronatus. Other palaeosubspecies are known, including those belonging to C. elaphus rianensis from the Middle Pleistocene of Italy, C. elaphus siciliae from the late Middle and Late Pleistocene of Sicily.
The International Union for Conservation of Nature originally listed nine subspecies of red deer (Cervus elaphus): three as endangered, one as vulnerable, one as near threatened, and four without enough data to give a category (Data Deficient). The species as a whole, however, is listed as least concern. However, this was based on the traditional classification of red deer as one species (Cervus elaphus), including the wapiti. The common red deer is also known as simply red deer.
Selected members of the red deer species group are listed in the table below. Of the ones listed, C. e. hippelaphus and C. e. scoticus may be junior synonyms.
Behaviour
Mature red deer (C. elaphus) usually stay in single-sex groups for most of the year. During the mating season, called the rut, mature stags compete for the attentions of the hinds and will then try to defend the hinds they attract. Rival stags challenge opponents by belling and walking in parallel. This allows combatants to assess each other's antlers, body size and fighting prowess. If neither stag backs down, a clash of antlers can occur, and stags sometimes sustain serious injuries. Red deer are among the mammals exhibiting homosexual behavior.
Dominant stags urinate on themselves and follow groups of hinds during the rut, from August into early winter. The stags may have as many as 20 hinds to keep from other, less attractive males. Only mature stags hold harems (groups of hinds), and breeding success peaks at about eight years of age. Stags two to four years old rarely hold harems and spend most of the rut on the periphery of larger harems, as do stags over 11 years old. Young and old stags that do acquire a harem hold it later in the breeding season than those stags in their prime. Harem-holding stags rarely feed and lose up to 20% of their body weight. Stags that enter the rut in poor condition are less likely to make it through to the peak conception period.
Male European red deer have a distinctive roar during the rut, which is an adaptation to forested environments, in contrast to male American elk stags which "bugle" during the rut in adaptation to open environments. The male deer roars to keep his harem of females together. The females are initially attracted to those males that both roar most often and have the loudest roar call. Males also use the roar call when competing with other males for females during the rut, and along with other forms of posturing and antler fights, is a method used by the males to establish dominance. Roaring is most common during the early dawn and late evening, which is also when the crepuscular deer are most active in general.
Breeding, gestation and lifespan
Female red deer reach sexual maturity at 2 years of age. Red deer mating patterns usually involve a dozen or more mating attempts before the first successful one. There may be several more matings before the stag will seek out another mate in his harem. Females in their second autumn can produce one or very rarely two offspring per year. The gestation period is 240 to 262 days, and the offspring weigh about . After two weeks, calves are able to join the herd and are fully weaned after two months. The offspring will remain with their mothers for almost one full year, leaving around the time the next season's offspring are produced. The gestation period is the same for all subspecies.
All red deer calves are born spotted, as is common with many deer species, and lose their spots by the end of summer. However, as in many species of Old World deer, some adults do retain a few spots on the backs of their summer coats.
Red deer live over 20 years in captivity and in the wild they live 10 to 13 years, though some subspecies with less predation pressure average 15 years.
Protection from predators
Male red deer retain their antlers for more than half the year, and are less gregarious and less likely to group with other males when they have antlers. The antlers provide self-defence, as does a strong front-leg kicking action performed by both sexes when attacked. Once the antlers are shed, stags tend to form bachelor groups which allow them to cooperatively work together. Herds tend to have one or more members watching for potential danger, while the remaining members eat and rest.
After the rut, females form large herds of up to 50 individuals. The newborn calves are kept close to the hinds by a series of vocalizations between the two, and larger nurseries have an ongoing and constant chatter during the daytime hours. When approached by predators, the largest and most robust females may make a stand, using their front legs to kick at their attackers. Guttural grunts and posturing is used with all but the most determined of predators with great effectiveness. Aside from humans and domestic dogs, the grey wolf is probably the most dangerous predator European red deer encounter. Occasionally, the brown bear will prey on European red deer.
Red deer in folklore and art
Red deer are widely depicted in cave art found throughout European caves, with some of the artwork dating from as early as 40,000 years ago, during the Upper Paleolithic. Siberian cave art from the Neolithic of 7,000 years ago has abundant depictions of red deer, including what can be described as spiritual artwork, indicating the importance of this mammal to the peoples of that region (Note: these animals were most likely wapiti (C. canadensis) in Siberia, not red deer).
Red deer are also often depicted on Pictish stones (circa 550–850 AD), from the early medieval period in Scotland, usually as prey animals for human or animal predators. In medieval hunting, the red deer was the most prestigious quarry, especially the mature stag, which in England was called a hart.
Red deer products
Red deer are held in captivity for a variety of reasons. The meat of the deer, called venison, was until recently restricted in the United Kingdom to those with connections to the aristocratic or poaching communities, and a licence was needed to sell it legally, but it is now widely available in supermarkets, especially in the autumn. The Queen followed the custom of offering large pieces of venison to members of the Cabinet of the United Kingdom and others. Some estates in the Scottish Highlands still sell deer-stalking accompanied by a gillie in the traditional way, on unfenced land, while others operate more like farms for venison. Venison is widely considered to be both flavourful and nutritious. It is higher in protein and lower in fat than either beef or chicken.
The red deer can produce of antler velvet annually. On ranches in New Zealand, China, Siberia, and elsewhere, this velvet is collected and sold to markets in East Asia, where it is used for holistic medicines, with South Korea being the primary consumer. In Russia, a medication produced from antler velvet is sold under the brand name Pantokrin (; ).
The antlers themselves are also believed by East Asians to have medicinal purposes and are often ground up and used in small quantities.
Historically, related deer species such as Central Asian red deer, wapiti, Thorold's deer, and sika deer have been reared on deer farms in Central and Eastern Asia by Han Chinese, Turkic peoples, Tungusic peoples, Mongolians, and Koreans. In modern times, western countries such as New Zealand and United States have taken to farming European red deer for similar purposes.
Deer hair products are also used in the fly fishing industry, being used to tie flies.
Deer antlers are also used for decorative purposes and have been used for artwork, furniture and other novelty items. Deer antlers were and still are the source material for horn furniture. Already in the 15th century trophies of case were used for clothes hook, storage racks and chandeliers, the so-called Lusterweibchen. In the 19th century the European nobility discovered red deer antlers as perfect decorations for their manors and hunting castles. This fashion trend splashes over to upper- and middle-class households in the mid of the 19th century.
At the increasingly popular World Expositions, producers of horn furniture, mainly in Germany, Austria and the United States, such as and Friedrich Wenzel, showed their horn furniture and a kind of series manufacturing began. In recent times deer antler home decors can be found in home styling magazines.
Gallery
| Biology and health sciences | Artiodactyla | null |
780166 | https://en.wikipedia.org/wiki/Menstrual%20pad | Menstrual pad | A menstrual pad is an absorbent item worn in the underwear when menstruating, bleeding after giving birth, recovering from gynecologic surgery, experiencing a miscarriage or abortion, or in any other situation where it is necessary to absorb a flow of blood from the vagina. A menstrual pad is a type of menstrual hygiene product that is worn externally, unlike tampons and menstrual cups, which are worn inside the vagina. Pads are generally changed by being stripped off the pants and panties, taking out the old pad, sticking the new one on the inside of the panties and pulling them back on. Pads are recommended to be changed every 3–4 hours to avoid certain bacteria that can fester in blood; this time also may differ depending on the kind worn, flow, and the time it is worn.
Menstrual pads are made from a range of materials, differing depending on style, country of origin, and brand. The pads are not the same as incontinence pads, which generally have higher absorbency and are worn by those who have urinary incontinence problems. Although menstrual pads are not made for this use, some use them for this purpose.
Types
Disposable
Although producers are generally reluctant to reveal the exact composition of their products, the main materials of disposable menstrual pads are usually bleached rayon (cellulose made from wood pulp), cotton and plastics. In addition, fragrance and antibacterial agents can be included. The plastic parts are the backsheet and polymer powder as an additional powerful absorbent (superabsorbent polymers) that turns into a gel when moistened. Procter & Gamble advertise a proprietary material called Infinicel as the core of their napkins.
In general, the layering is as follows: "an absorbent core material placed between a flexible liquid-pervious topsheet and a liquid-impervious plastic backsheet that has an adhesive on the outside for attaching the napkin to an undergarment". As is the case with disposable tampons and diapers, recycling is difficult and rarely done due to cost reasons, although proof-of-principle solutions appear to exist. When not dumped in a landfill where the non-biodegradable parts may persist for thousands of years, conventional hygiene products can at best be "thermally recycled" (incinerated).
There are several different types of disposable menstrual pads:
Panty liner: Designed to absorb daily vaginal discharge, light menstrual flow, spotting, slight urinary incontinence, or as a backup for tampon or menstrual cup use.
Ultra-thin: A very compact which may be as absorbent as a Regular or Maxi/Super pad but with less bulk.
Regular: A middle range absorbency pad.
Maxi/Super: A larger absorbency pad, useful for the start of the menstrual cycle when menstruation is often heaviest.
Overnight: A longer pad to allow for more protection while the wearer is lying down, with an absorbency suitable for overnight use.
Maternity: These are usually slightly longer than a maxi/Super pad and are designed to be worn to absorb lochia (bleeding that occurs after childbirth) and also can absorb urine.
The shape, absorbency and lengths may vary depending on manufacturer, but usually range from the short slender panty liner to the larger and longer overnight. Long pads are offered for extra protection or for larger people whose undergarments might not be completely protected by regular length pads, and also for overnight use.
Other options are often offered in a manufacturer's line of pads, such as wings or tabs that wrap around the sides of underwear to add additional leak protection and help secure the pad in place. Deodorant is also added to some pads, which is designed to cover menstrual odor with a light fragrance. There are even panty liners specifically designed to be worn with a thong/G-string.
Reusable (cloth)
Some people use a washable or reusable cloth menstrual pad. These are made from a number of types of fabric—most often cotton flannel, bamboo, or hemp (which is highly absorbent and not as bulky as cotton). Most styles have wings that secure around the underpants, but some are just held in place (without wings) between the body and the underpants. Some (particularly the older styles) are available in belted styles. Cloth menstrual pads made a comeback around the 1970s, with their popularity increasing in the late 1980s and early 1990s. Reasons people choose to switch to cloth menstrual pads include comfort, savings over time, environmental impact and health reasons.
Washable menstrual pads do not need to be disposed of after use and therefore offer a more economical alternative. Reusable menstrual pads can be found on a number of websites, or are made at home (instructions are available online). They have become a popular alternative because they are allergen- and perfume-free, and can be more comfortable for people who experience irritation from using disposable pads.
There are many styles of cloth menstrual pads available today, ranging from panty liners to overnight pads. Popular styles of cloth menstrual pads include all-in-one, or AIO pads, in which the absorbent layer is sewn inside the pad, "inserts on top" style pads, which have absorbent layers that can be secured on top of the pad as needed, envelope or pocket style pads, which have absorbent layers that can be inserted inside the pad as needed, and a foldable style, in which the pad folds around the absorbent layers. Cloth menstrual pads can also have a waterproof lining, which provides more leak protection but may also be less breathable than those with no lining.
Uses
Menstrual pads are worn to absorb menstrual discharge, thereby protecting clothing and furnishings. They are usually individually wrapped so they are easier and more discreet to carry in a purse or bag. This wrapper may be used to wrap the soiled pads before disposing of them in appropriate receptacles. Some people prefer to wrap the pads with toilet paper instead of (or as well as) using the wrapper, which, often being made of slick plastic with a small tape tab, may not adequately stick. Menstrual pads of any type should not be flushed down the toilet as they can cause blockages. In developed countries, public toilets almost always include a receptacle in which to place soiled pads. In first aid, they make excellent dressings for heavy bleeding due to their high absorbency if gauze is unavailable or inadequate.
Many people who experience urinary incontinence use menstrual pads to manage bladder leaks. However, since menstrual pads are designed to absorb menstrual flow, they are not as effective in absorbing urinary leaks; incontinence pads are designed for this purpose.
If someone does not have menstrual pads on hand while menstruating, they might use toilet paper as a temporary substitution.
History
Through the ages societies have used different forms of menstrual protection. Menstrual pads have been mentioned as early as the 10th century, in the Suda, where Hypatia, who lived in the 4th century AD, was said to have thrown one of her used menstrual rags at an admirer in an attempt to discourage him. In ancient Egypt, women used softened papyrus, a grass-like plant, to absorb their menstrual blood. Before commercially available menstrual hygiene products, most women used pieces of cloth to absorb their menstrual flow. "On the rag" is a term that originally referred to menstrual rags, but its modern usage is as a menstrual euphemism.
Until disposable menstrual pads were created, cloth or reusable pads were widely used to collect menstrual blood. People often used a variety of home-made menstrual pads which they crafted from various fabrics, or other absorbent materials, to collect menstrual blood. Even after disposable pads were commercially available, for several years they were too expensive for many families to afford. When they could be afforded, people were allowed to place money in a box so that they would not have to speak to the clerk and take a box of Kotex pads from the counter themselves. It took several years for disposable menstrual pads to become commonplace. For easier use, inventor Carolyn R. Mobley patented the tab construction for a sanitary napkin. Disposables are now used nearly exclusively in most of the industrialized world.
The first of the disposable pads were generally in the form of a cotton wool or similar fibrous rectangle covered with an absorbent liner. The liner ends were extended front and back so as to fit through loops in a special girdle or belt worn beneath undergarments. This design was notorious for slipping either forward or back of the intended position.
Disposable menstrual pads grew from a Benjamin Franklin invention created to help stop wounded soldiers from bleeding, but appear to have been first commercially available from around 1880 with Thomas and William Southall's pad. The first commercially available American disposable napkins were Lister's Towels created by Johnson & Johnson in 1888. Disposable pads had their start with nurses using their wood pulp bandages to absorb their menstrual flow, creating a pad that was made from easily obtainable materials and inexpensive enough to throw away after use. Kotex's first advertisement for products made with this wood pulp (cellucotton) appeared in January 1921. Johnson & Johnson introduced Modess Sanitary Napkins in 1926, researched by Lillian Gilbreth. Lilian Gilbreth's market research report published in 1927 gives valuable information about American's experiences of menstruation in 1920s. The surveys she conducted on over 1000 people with periods reveal that "the most significant fact concerning the marketing angle of the sanitary napkin [...] is that it is availability that sells a napkin." Several of the first disposable pad manufacturers were also manufacturers of bandages, which could give an indication of what these products were like.
In 1956, Mary Kenner obtained a patent for an adjustable sanitary belt with an inbuilt, moisture-proof napkin pocket. However, the company that first showed interest in her invention rejected it after they discovered that she was African American.
Later an adhesive strip was placed on the bottom of sanitary pads for attachment to the crotch of the panties, and this became a favoured method. The belted sanitary napkin quickly disappeared during the early 1980s.
The ergonomic design and materials used to make pads also changed through the 1980s to today. With earlier materials not being as absorbent and effective, and early pads being up to two centimetres thick, leaks were a major problem. Some variations introduced were quilting of the lining, adding "wings" and reducing the thickness of the pad by utilising products such as sphagnum and polyacrylate superabsorbent gels derived from petroleum. The materials used to manufacture most pads are derived from the petroleum industry and forestry. The absorbent core, made from chlorine bleached wood pulp, could be reduced to make slimmer products with the addition of polyacrylate gels which sucks up the liquid quickly and holds it in a suspension under pressure. The remaining materials are mostly derived from the petroleum industry, the cover stock used is polypropylene non woven, with the leakproof barrier made from polyethylene film.
Society and culture
Pads, especially reusable ones, may be visible on full body scanners.
Developing countries
In developing countries, makeshift pads are still used to collect menstrual blood as they are cheaper. Rags, soil, and mud are also reportedly used for collecting menstrual flow by people who cannot afford the more expensive disposable pads or tampons.
In order to meet the need for achieving an inexpensive solution to reduce unsanitary and unhygienic practices in countries like India, Arunachalam Muruganantham from rural Coimbatore in the southern state of Tamil Nadu, India, developed and patented a machine which could manufacture low-cost sanitary pads for less than a third of the former cost. The Bill and Melinda Gates Foundation awarded Nairobi-based ZanaAfrica a grant of US$3 million. ZanaAfrica explores creative approaches to menstrual health education for adolescent girls.
Health effects
In 2024 it was reported that a brand of menstrual pad was found to contain the toxic chemical PFOA.
| Biology and health sciences | Hygiene products | Health |
780489 | https://en.wikipedia.org/wiki/Thylacoleo | Thylacoleo | Thylacoleo ("pouch lion") is an extinct genus of carnivorous marsupials that lived in Australia from the late Pliocene to the Late Pleistocene (until around 40,000 years ago), often known as marsupial lions. They were the largest and last members of the family Thylacoleonidae, occupying the position of apex predator within Australian ecosystems. The largest and last species, Thylacoleo carnifex, approached the weight of a lioness. The estimated average weight for the species ranges from .
Taxonomy
The first Thylacoleo fossil findings, discovered by Thomas Mitchell were found in the 1830s in the Wellington Valley of New South Wales, though not recognised as such at the time. The generic holotype, consisting of broken teeth, jaws, and a skull was discovered by a pastoralist, William Avery, near Lake Colungolac from which the species Thylacoleo carnifex was described by Richard Owen. It was not until 1966 that the first nearly-complete skeleton was found. The only pieces missing were a foot and the tail. Currently, the Nullarbor Plain of West Australia remains to be the greatest finding site. These fossils now reside at the Australian Museum.
The genus was first published in 1859, erected to describe the type species Thylacoleo carnifex. The new taxon was established in examination of fossil specimens provided to Richard Owen. The familial alliance takes its name from this description, the so-called marsupial lions of Thylacoleonidae.
The colloquial name "marsupial lion" alludes to the genus name, which was named after its superficial resemblance to the placental lion and its ecological niche as a large predator.
Genus: Thylacoleo (Thylacopardus) – Australia's marsupial lions, that lived from about 2 million years ago, during the Late Pliocene Epoch and became extinct about 40,000 years ago, during the Late Pleistocene Epoch. Three species are known:
Thylacoleo carnifex The holotype cranium was collected from Lake Colongulac in 1843 by pastoralist William Adeney. A partial rostrum collected by Adeney in 1876 from the same locality would later be found to belong to the same individual. It was not until 1966 that the first nearly-complete skeleton was found.
Thylacoleo crassidentatus lived during the Pliocene, around 5 million years ago, and was about the size of a large dog. Its fossils have been found in southeastern Queensland.
Thylacoleo hilli lived during the Pliocene and was half the size of T. crassidentatus. It is the oldest member of the genus.
Fossils of other representatives of Thylacoleonidae, such as Microleo and Wakaleo, date back to the Late Oligocene Epoch, some 24 million years ago.
T. hilli was described by Neville Pledge in a study published in the records of the South Australia Museum in 1977. The holotype is a third premolar, discovered at a cave in Curramulka in South Australia, exhibiting the carnivorous characteristics of the genus and around half the size of T. carnifex. This tooth was collected by Alan Hill, a speleologist and founding member of the Cave Exploration Group of South Australia, while examining a site known as the "Town Cave" in 1956; the specific epithet hilli honours the collector of the first specimen. Material found amidst the fauna at Bow River in New South Wales, dated to the early Pliocene, was also referred to the species in 1982.
A fragment of an incisor, unworn and only diagnosable to the genus, was located at a site in Curramulka, close to the Town Cave site, and referred to the species for the apparent correlation in size when compared to the better known T. carnifex.
The marsupial lion is classified in the order Diprotodontia along with many other well-known marsupials such as kangaroos, possums, and the koala. It is further classified in its own family, the Thylacoleonidae, of which three genera and 11 species are recognised, all extinct. The term marsupial lion (lower case) is often applied to other members of this family. Distinct possum-like characteristics led Thylacoleo to be regarded as members of Phalangeroidea for a few decades. Though a few authors continued to hint at phalangeroid affinities for thylacoleonids as recently as the 1990s, cranial and other characters have generally led to their inclusion within vombatiformes, and as stem-members of the wombat lineage. Marsupial lions and other ecologically and morphologically diverse vombatiforms were once represented by over 60 species of carnivorous, herbivorous, terrestrial and arboreal forms ranging in size from 3 kg to 2.5 tonnes. Only two families represented by four herbivorous species (koalas and three species of wombat) have survived into modern times and are considered the marsupial lion's closest living relatives.
Evolution
The ancestors of thylacoleonids are believed to have been herbivores, something unusual for carnivores. They are members of the Vombatiformes, an almost entirely herbivorous order of marsupials, the only extant representatives of which are koalas and wombats, as well as extinct members such as the diprotodontids and palorchestids. The group first appeared in the Late Oligocene. The earliest thylacoleonids like Microleo were small possum-like animals, with the group increasing in size during the Miocene, with representatives like the leopard-sized Wakaleo. The genus Thylacoleo first appeared during the Pliocene, and represented the only extant genus of the family from that time until the end of the Pleistocene. The youngest representative of Thylacoleo and the thylacoleonids, T. carnifex, is the largest known member of the family. The earliest thylacoleonids are thought to have been arboreal (tree dwelling) animals, while Thylacoleo is thought to be terrestrial with some climbing capabilities.
Description
T. carnifex is the largest carnivorous mammal known to have ever existed in Australia, and one of the largest metatherian carnivores known (comparable to Thylacosmilus and Borhyaena species, but smaller than Proborhyaena gigantea). Individuals ranged up to around high at the shoulder and about from head to tail. Measurements taken from a number of specimens show they averaged in weight, although individuals as large as might not have been uncommon, and the largest weight was of . This would make it comparable to female lions and female tigers in general size. Estimates of the size of T. carnifex based on dental remains are typically dubious, in contrast to estimates based on proximal limb bone circumference.
Skull
Like other thylacoleonids, Thylacoleo had blade-like third premolar teeth in the upper and lower jaws, that functioned as the carnassial teeth, with these teeth being present much further forwards in the jaw than in other mammals. Compared to earlier thyacoleonids, the third premolars were considerably enlarged. Thylacoleo also had a proportionally large pair of first incisors in the upper and lower jaws, which functioned analogously to other carnivores canine teeth. They also had true canines but they served little purpose as they were stubby and not very sharp. Compared to earlier thylacoleonids, the number of molar teeth was reduced.
Pound for pound, T. carnifex had the strongest bite of any mammal species, living or extinct; a T. carnifex weighing had a bite comparable to that of a 250 kg African lion, and research suggests that Thylacoleo could hunt and take prey much larger than itself. Larger animals that were likely prey include Diprotodon spp. and giant kangaroos. It seems improbable that Thylacoleo could achieve as high a bite force as a modern-day lion; however, this might have been possible when taking into consideration the size of its brain and skull. Carnivores usually have rather large brains when compared to herbivorous marsupials, which lessens the amount of bone that can be devoted to enhancing bite force. Thylacoleo however, is thought to have had substantially stronger muscle attachments and therefore a smaller brain. Some later studies questioned the ability of the canine teeth to deliver a killing bite.
Using 3D modeling based on X-ray computed tomography scans, marsupial lions were found to be unable to use the prolonged, suffocating bite typical of living big cats. They instead had an extremely efficient and unique bite; the incisors would have been used to stab at and pierce the flesh of their prey while the more specialised carnassials crushed the windpipe, severed the spinal cord, and lacerated the major blood vessels such as the carotid artery and jugular vein. Compared to an African lion which may take 15 minutes to kill a large catch, the marsupial lion could kill a large animal in less than a minute. The skull was so specialized for big game that it was very inefficient at catching smaller animals, which possibly contributed to its extinction.
Postcranium
Thylacoleo had highly mobile and powerful forelimbs used to grapple prey, with each manus having a single very large retractable hooked claw set on large semi-opposable thumbs, which are suggested to have been used deal a killing blow.
The hind feet had four functional toes, the first digit being much reduced in size, but possessing a roughened pad similar to that of possums, which may have assisted with climbing. The discovery in 2005 of a specimen which included complete hind feet provided evidence that the marsupial lion exhibited syndactyly (fused second and third toes) like other diprotodonts.
Its strong forelimbs and retracting claws mean that Thylacoleo possibly climbed trees and perhaps carried carcasses to keep the kill for itself (similar to the leopard today). The climbing ability would have also helped them climb out of caves, which could therefore have been used as dens to rear their young. Specialised tail bones called chevrons strengthened the tail, likely allowing the animal to use it to prop itself up while rearing on its hind legs, which may have been done when climbing or attacking prey.
The lumbar region is relatively rigid and straight, and suggests that the lower back was relatively inflexible.
Ecology
When Thylacoleo was first described by Richard Owen, he considered it to be a carnivore, based on the morphology of its skull and teeth. However other anatomists, such as William Henry Flower disagreed. Flower was the first to place Thylacoleo with the Diprotodonts, noting its skull and teeth to be laid out more like those of the koala and the wombat, and suggested that it was more likely a herbivore. Owen did not disagree with Flower's placement of Thylacoleo with the Diprotodonts, but still maintained that it was a carnivore, despite its herbivorous ancestry. Owen found little support in his lifetime, despite the pointing out of Thylacoleo'''s retractable claws, something only found in mammalian carnivores, and its lack of any ability to chew plant material. In 1911, a study by Spencer and Walcott claimed that certain marks on the bones of megafauna had been made by Thylacoleo, but according to Horton (1979) they were not sufficiently rigorous, resulting in their arguments being strongly challenged by later scholars, such as Anderson (1929), and later Gill (1951, 1952, 1954), thereby leaving the issue unresolved.
Besides the most common hypothesis that it was an active predator, a variety of other theories existed in the late 19th to early 20th centuries as to the diet and feeding of Thylacoleo, with hypotheses of it being a scavenger filling the ecological niche of hyenas, being a specialist of crocodile eggs, or even a melon-eater. As late as 1954, doubts were still being raised as to whether it was actually a hypercarnivore.
In 1981, another paper was published arguing that certain cuts to bones of large marsupials had been caused by Thylacoleo. This paper by Horton and Wright was able to counter earlier arguments that such marks were the result of humans, largely by pointing out the presence of similar marks on the opposite side of many bones. They concluded that humans were extremely unlikely to have made the marks in question, but that if so "they had set out to produce only marks consistent with what Thylacoleo would produce". Since then, the academic consensus has emerged that Thylacoleo was a predator and a hypercarnivore.
The marsupial lion's limb proportions and muscle mass distribution indicate that, although it was a powerful animal, it was not a particularly fast runner. Paleontologists conjecture that it was an ambush predator, possibly using leaping. Incisions on bones of the extinct large kangaroo Macropus titan, and the general morphology of Thylacoleo suggests that it fed in a similar manner to modern cheetahs, by using their sharp teeth to slice open the ribcage of their prey, thereby accessing the internal organs. They may have killed by using their front claws as either stabbing weapons or as a way to grab their prey with strangulation or suffocation.
Like many predators, it was probably also an opportunistic scavenger, feeding on carrion and driving off less powerful predators from their kills. It also may have shared behaviours exhibited by recent diprotodont marsupials such as kangaroos, like digging shallow holes under trees to reduce body temperature during the day.
Trace fossils in the form of claw marks and bones from caves in Western Australia analyzed by Gavin Prideaux et al. indicate marsupial lions could also climb rock faces, and likely reared their young in such caves as a way of protecting them from potential predators.
Analysis of finds on the Nullabor Plain suggests that Thylacoleo carnifex inhabited open, arid environments similar to those found across much of Australia today.
A 1985 study suggested that Thylacoleo carnifex was an apex predator that primarily fed on large bodied prey, which may have included the large kangaroos Sthenurus, Procoptodon, Protemnodon, Macropus and Ostphranter as well as possibly the largest Australian marsupial, the rhinoceros-sized vombatoid Diprotodon.
Extinction Thylacoleo'' is thought to have become extinct around 40,000 years ago as part of the Late Pleistocene megafauna extinctions, essentially simultanteously with the vast majority of Australian megafauna. It has been contested as to the relative importance of climatic change vs the impact of recently arrived Indigenous Australians (who arrived in Australia around 50-60,000 years ago) in the extinctions. There is limited evidence of human interaction with extinct megafauna in Australia.
| Biology and health sciences | Diprotodontia | Animals |
780629 | https://en.wikipedia.org/wiki/Non-volatile%20memory | Non-volatile memory | Non-volatile memory (NVM) or non-volatile storage is a type of computer memory that can retain stored information even after power is removed. In contrast, volatile memory needs constant power in order to retain data.
Non-volatile memory typically refers to storage in memory chips, which store data in floating-gate memory cells consisting of floating-gate MOSFETs (metal–oxide–semiconductor field-effect transistors), including flash memory storage such as NAND flash and solid-state drives (SSD).
Other examples of non-volatile memory include read-only memory (ROM), EPROM (erasable programmable ROM) and EEPROM (electrically erasable programmable ROM), ferroelectric RAM, most types of computer data storage devices (e.g. disk storage, hard disk drives, optical discs, floppy disks, and magnetic tape), and early computer storage methods such as punched tape and cards.
Overview
Non-volatile memory is typically used for the task of secondary storage or long-term persistent storage. The most widely used form of primary storage today is a volatile form of random access memory (RAM), meaning that when the computer is shut down, anything contained in RAM is lost. However, most forms of non-volatile memory have limitations that make them unsuitable for use as primary storage. Typically, non-volatile memory costs more, provides lower performance, or has a limited lifetime compared to volatile random access memory.
Non-volatile data storage can be categorized into electrically addressed systems, for example, flash memory, and read-only memory) and mechanically addressed systems (hard disks, optical discs, magnetic tape, holographic memory, and such). Generally speaking, electrically addressed systems are expensive, and have limited capacity, but are fast, whereas mechanically addressed systems cost less per bit, but are slower.
Electrically addressed
Electrically addressed semiconductor non-volatile memories can be categorized according to their write mechanism.
Read-only and read-mostly devices
Mask ROMs are factory programmable only and typically used for large-volume products which are not required to be updated after the memory device is manufactured.
Programmable read-only memory (PROM) can be altered once after the memory device is manufactured using a PROM programmer. Programming is often done before the device is installed in its target system, typically an embedded system. The programming is permanent, and further changes require the replacement of the device. Data is stored by physically altering (burning) storage sites in the device.
An EPROM is an erasable ROM that can be changed more than once. However, writing new data to an EPROM requires a special programmer circuit. EPROMs have a quartz window that allows them to be erased with ultraviolet light, but the whole device is cleared at one time. A one-time programmable (OTP) device may be implemented using an EPROM chip without the quartz window; this is less costly to manufacture. An electrically erasable programmable read-only memory EEPROM uses voltage to erase memory. These erasable memory devices require a significant amount of time to erase data and write new data; they are not usually configured to be programmed by the processor of the target system. Data is stored using floating-gate transistors, which require special operating voltages to trap or release electric charge on an insulated control gate to store information.
Flash memory
Flash memory is a solid-state chip that maintains stored data without any external power source. It is a close relative to the EEPROM; it differs in that erase operations must be done on a block basis, and its capacity is substantially larger than that of an EEPROM. Flash memory devices use two different technologies—NOR and NAND—to map data. NOR flash provides high-speed random access, reading and writing data in specific memory locations; it can retrieve as little as a single byte. NAND flash reads and writes sequentially at high speed, handling data in blocks. However, it is slower on reading when compared to NOR. NAND flash reads faster than it writes, quickly transferring whole pages of data. Less expensive than NOR flash at high densities, NAND technology offers higher capacity for the same-size silicon.
Ferroelectric RAM (F-RAM)
Ferroelectric RAM (FeRAM, F-RAM or FRAM) is a form of random-access memory similar in construction to DRAM, both use a capacitor and transistor but instead of using a simple dielectric layer the capacitor, an F-RAM cell contains a thin ferroelectric film of lead zirconate titanate , commonly referred to as PZT. The Zr/Ti atoms in the PZT change polarity in an electric field, thereby producing a binary switch. Due to the PZT crystal maintaining polarity, F-RAM retains its data memory when power is shut off or interrupted.
Due to this crystal structure and how it is influenced, F-RAM offers distinct properties from other nonvolatile memory options, including extremely high, although not infinite, endurance (exceeding 1016 read/write cycles for 3.3 V devices), ultra-low power consumption (since F-RAM does not require a charge pump like other non-volatile memories), single-cycle write speeds, and gamma radiation tolerance.
Magnetoresistive RAM (MRAM)
Magnetoresistive RAM stores data in magnetic storage elements called magnetic tunnel junctions (MTJs). The first generation of MRAM, such as Everspin Technologies' 4 Mbit, utilized field-induced writing. The second generation is developed mainly through two approaches: Thermal-assisted switching (TAS) which is being developed by Crocus Technology, and Spin-transfer torque (STT) which Crocus, Hynix, IBM, and several other companies are developing.
Phase-change Memory (PCM)
Phase-change memory stores data in chalcogenide glass, which can reversibly change the phase between the amorphous and the crystalline state, accomplished by heating and cooling the glass. The crystalline state has low resistance, and the amorphous phase has high resistance, which allows currents to be switched ON and OFF to represent digital 1 and 0 states.
FeFET memory
FeFET memory uses a transistor with ferroelectric material to permanently retain state.
RRAM memory
RRAM (ReRAM) works by changing the resistance across a dielectric solid-state material often referred to as a memristor. ReRAM involves generating defects in a thin oxide layer, known as oxygen vacancies (oxide bond locations where the oxygen has been removed), which can subsequently charge and drift under an electric field. The motion of oxygen ions and vacancies in the oxide would be analogous to the motion of electrons and holes in a semiconductor.
Although ReRAM was initially seen as a replacement technology for flash memory, the cost and performance benefits of ReRAM have not been enough for companies to proceed with the replacement. Apparently, a broad range of materials can be used for ReRAM. However, the discovery that the popular high-κ gate dielectric HfO2 can be used as a low-voltage ReRAM has encouraged researchers to investigate more possibilities.
Mechanically addressed systems
Mechanically addressed systems use a recording head to read and write on a designated storage medium. Since the access time depends on the physical location of the data on the device, mechanically addressed systems may be sequential access. For example, magnetic tape stores data as a sequence of bits on a long tape; transporting the tape past the recording head is required to access any part of the storage. Tape media can be removed from the drive and stored, giving indefinite capacity at the cost of the time required to retrieve a dismounted tape.
Hard disk drives use a rotating magnetic disk to store data; access time is longer than for semiconductor memory, but the cost per stored data bit is very low, and they provide random access to any location on the disk. Formerly, removable disk packs were common, allowing storage capacity to be expanded. Optical discs store data by altering a pigment layer on a plastic disk and are similarly random access. Read-only and read-write versions are available; removable media again allows indefinite expansion, and some automated systems (e.g. optical jukebox) were used to retrieve and mount disks under direct program control.
Domain-wall memory (DWM) stores data in a magnetic tunnel junctions (MTJs), which works by controlling domain wall (DW) motion in ferromagnetic nanowires.
Organic
Thinfilm produces rewriteable non-volatile organic ferroelectric memory based on ferroelectric polymers. Thinfilm successfully demonstrated roll-to-roll printed memories in 2009. In Thinfilm's organic memory the ferroelectric polymer is sandwiched between two sets of electrodes in a passive matrix. Each crossing of metal lines is a ferroelectric capacitor and defines a memory cell.
Non-volatile main memory
Non-volatile main memory (NVMM) is primary storage with non-volatile attributes. This application of non-volatile memory presents security challenges. NVDIMM is one example of the non-volatile main memory.
| Technology | Non-volatile memory | null |
780960 | https://en.wikipedia.org/wiki/Software%20maintenance | Software maintenance | Software maintenance is the modification of software after delivery.
Software maintenance is often considered lower skilled and less rewarding than new development. As such, it is a common target for outsourcing or offshoring. Usually, the team developing the software is different from those who will be maintaining it. The developers lack an incentive to write the code to be easily maintained. Software is often delivered incomplete and almost always contains some bugs that the maintenance team must fix. Software maintenance often initially includes the development of new functionality, but as the product nears the end of its lifespan, maintenance is reduced to the bare minimum and then cut off entirely before the product is withdrawn.
Each maintenance cycle begins with a change request typically originating from an end user. That request is evaluated and if it is decided to implement it, the programmer studies the existing code to understand how it works before implementing the change. Testing to make sure the existing functionality is retained and the desired new functionality is added often comprises the majority of the maintenance cost.
Software maintenance is not as well studied as other phases of the software life cycle, despite comprising the majority of costs. Understanding has not changed significantly since the 1980s. Software maintenance can be categorized into several types depending on whether it is preventative or reactive and whether it is seeking to add functionality or preserve existing functionality, the latter typically in the face of a changed environment.
History
In the early 1970s, companies began to separate out software maintenance with its own team of engineers to free up software development teams from support tasks. In 1972, R. G. Canning published "The Maintenance 'Iceberg, in which he contended that software maintenance was an extension of software development with an additional input: the existing system. The discipline of software maintenance has changed little since then. One twenty-first century innovation has been companies deliberately releasing incomplete software and planning to finish it post-release. This type of change, and others that expand functionality, is often called software evolution instead of maintenance.
Software life cycle
Despite testing and quality assurance, virtually all software contains bugs where the system does not work as intended. Post-release maintenance is necessary to remediate these bugs when they are found. Most software is a combination of pre-existing commercial off-the-shelf (COTS) and open-source software components with custom-written code. COTS and open-source software is typically updated over time, which can reduce the maintenance burden, but the modifications to these software components will need to be adjusted for in the final product. Unlike software development, which is focused on meeting specified requirements, software maintenance is driven by events—such as user requests or detection of a bug. Its main purpose is to preserve the usefulness of the software, usually in the face of changing requirements.
If conceived of as part of the software development life cycle, maintenance is the last and typically the longest phase of the cycle, comprising 80 to 90 percent of the lifecycle cost. Other models consider maintenance separate from software development, instead as part of the software maintenance life cycle (SMLC). SMLC models typically include understanding the code, modifying it, and revalidating it.
Transition from release to maintenance to end of the lifespan
Frequently, software is delivered in an incomplete state. Developers will test a product until running out of time or funding, because they face fewer consequences for an imperfect product than going over time or budget. The transition from the development to the maintenance team is often inefficient, without lists of known issues or validation tests, which the maintenance team will likely recreate. After release, members of the development team are likely to be reassigned or otherwise become unavailable. The maintenance team will require additional resources for the first year after release, both for technical support and fixing defects left over from development.
Initially, software may go through a period of enhancements after release. New features are added according to user feedback. At some point, the company may decide that it is no longer profitable to make functional improvements, and restrict support to bug fixing and emergency updates. Changes become increasingly difficult and expensive due to lack of expertise or decaying architecture due to software aging. After a product is no longer maintained, and does not receive even this limited level of updating, some vendors will seek to extract revenue from the software as long as possible, even though the product is likely to become increasingly avoided. Eventually, the software will be withdrawn from the market, although it may remain in use. During this process, the software becomes a legacy system.
Change cycle
The first step in the change cycle is receiving a change request from a customer and analyzing it to confirm the problem and decide whether to implement the change. This may require input from multiple departments; for example, the marketing team can help evaluate whether the change is expected to bring more business. Software development effort estimation is a difficult problem, including for maintenance change requests, but the request is likely to be declined if it is too expensive or infeasible. If it is decided to implement the request, it can be assigned to a scheduled release and implemented. Although agile methodology does not have a maintenance phase, the change cycle can be enacted as a scrum sprint.
Understanding existing code is an essential step before modifying it. The rate of understanding depends both on the code base as well as the skill of the programmer. Following coding conventions such as using clear function and variable names that correspond to their purpose makes understanding easier. Use of conditional loop statements only if the code could execute more than once, and eliminating code that will never execute can also increase understandability. Experienced programmers have an easier time understanding what the code does at a high level. Software visualization is sometimes used to speed up this process.
Modification to the code may take place in any way. On the one hand, it is common to haphazardly apply a quick fix without being granted enough time to update the code documentation. On the other hard structured iterative enhancement can begin by changing the top-level requirements document and propagating the change down to lower levels of the system. Modification often includes code refactoring (improving the structure without changing functionality) and restructuring (improving structure and functionality at the same time). Unlike commercial software, free and open source software change cycles are largely restricted to coding and testing, with minimal documentation. Open-source software projects instead rely on mailing lists and a large number of contributors to understand the code base and fix bugs efficiently.
An additional problem with maintenance is that nearly every change to code will introduce new bugs or unexpected ripple effects, which require another round of fixes. Testing can consume the majority of maintenance resource for safety-critical code, due to the need to revalidate the entire software if any changes are made. Revalidation may include code review, regression testing with a subset of unit tests, integration tests, and system tests. The goal of the testing is to verify that previous functionality is retained, and the new functionality has been added.
Categories of software maintenance
The key purpose of software maintenance is ensuring that the product continues to meet usability requirements. At times, this may mean extending the product's capabilities beyond what was initially envisioned.
According to the ISO/IEC 14764 specification, software maintenance can be classified into four types:
Corrective maintenance: modification of software to fix a bug or other failure to meet requirements, typically reported by an end user.
Preventive maintenance: forward-looking modification of software after delivery to ensure it continues to meet requirements or fix problems that have not manifested yet. This type of maintenance is performed especially on systems that are required to be highly safe or available. Software rejuvenation is one form of preventative maintenance to clean up state and prevent future problems.
Adaptive maintenance: modification of software performed after delivery to ensure its continuing usability in a changed or changing environment.
Perfective maintenance: enhancement of software after delivery to improve qualities such as user experience, processing efficiency, and maintainability. Perfective maintenance is necessary if other types of maintenance are carried out, because modification of an existing code base will otherwise increase complexity and cause the existing structure to deteriorate. Perfective maintenance may include rewriting documentation, code refactoring, and performance tuning.
According to some estimates, enhancement (the latter two categories) comprises some 80 percent of software maintenance.
Maintainability
Maintainability is the quality of software enabling it to be easily modified without breaking existing functionality. According to the ISO/IEC 14764 specification, activity to ensure software maintainability prior to release counts as part of software maintenance. Many software development organizations neglect maintainability, even though doing so will increase long-term costs. Technical debt is incurred when programmers, often out of laziness or urgency to meet a deadline, choose quick and dirty solutions rather than build maintainability into their code. A common cause is underestimates in software development effort estimation, leading to insufficient resources allocated to development. One important aspect is having a large amount of automated software tests that can detect if existing functionality is compromised by a change.
A challenge with maintainability is that many software engineering courses do not emphasize it, and give out one-and-done assignments that have clear and unchanging specifications. Software engineering courses do not cover systems as complex as occur in the real world. Development engineers who know that they will not be responsible for maintaining the software do not have an incentive to build in maintainability.
Workforce
Maintenance is often considered an unrewarding job for software engineers, who, if assigned to maintenance, were more likely to quit. It often pays less than a comparable job in software development. The task is often assigned to temporary workers or lesser-skilled staff, although maintenance engineers are also typically older than developers, partly because they must be familiar with outdated technologies. In 2008, around 900,000 of the 1.3 million software engineers and programmers working in the United States were doing maintenance.
Companies started separate teams for maintenance, which led to outsourcing this work to a different company, and by the turn of the twenty-first century, sometimes offshoring the work to another country—whether as part of the original company or a separate entity. The typical sources of outsourcing are developed countries such as the United States, the United Kingdom, Japan, and Australia, while destinations are usually lower-cost countries such as China, India, Russia, and Ireland. Reasons for offshoring include taking advantage of lower labor costs, enabling around-the-clock support, reducing time pressure on developers, and to move support closer to the market for the product. Downsides of offshoring include communication barriers in the form of such factors as time zone and organizational disjunction and cultural differences. Despite many employers considering maintenance lower-skilled work and the phase of software development most suited to offshoring, it requires close communication with the customer and rapid response, both of which are hampered by these communication difficulties.
Alternatives to maintenance
In software engineering, the term legacy system does not have a fixed meaning, but often refers to older systems which are large, difficult to modify, and also necessary for current business needs. Often legacy systems are written in obsolete programming languages, lack documentation, have a deteriorating structure after years of changes, and depend on experts to keep it operational. When dealing with these systems, at some point so much technical debt accumulates that maintenance is not practical or economical. Other choices include:
Freezing—do no more work on the legacy system. This option may be chosen if the vendor wants to continue to extract revenue as long as possible while avoiding maintenance costs.
Outsourcing functionality of the legacy system to a different company, especially if it is not considered a core business function.
Discarding the existing legacy system and redeveloping a new application from scratch to fulfill the same purpose as the legacy system. However, this approach is inefficient due to discarding a working system, and with this approach there is a danger that the new system will not fulfill changing business requirements.
Wrapping the legacy application in an abstraction layer to simplify outdated interfaces. The source code is not modified but the new interface allows a tried and tested component to be accessed by newer applications. This approach does not fix any of the issues with maintaining a legacy system. Databases, functions, and entire applications may be wrapped in this way.
Migrating the legacy system to a new platform, which can reduce the expense of new software development by reusing the implementation, design, specification, and requirements of the legacy system. Migration can take 5 to 10 years, but results in greater flexibility and long-term savings in software maintenance. As much as 80 percent of the expense is in testing; that is, ensuring that the new system has the same output as the old system. After the new system is finished, there needs to be a transition from the old system to the new system with minimum disruption to business functions.
Research
Despite taking up the lion's share of software development resources, maintenance is the least studied phase of software development. Much of the literature has focused on how to develop maintainable code from the outset, with less focus on motivating engineers to make maintainability a priority. , automated solutions for code refactoring to reduce maintenance effort are an active area of research, as is machine-learning enhanced maintainability assessment.
| Technology | Software development: General | null |
782310 | https://en.wikipedia.org/wiki/Geological%20hazard | Geological hazard | A geologic hazard or geohazard is an adverse geologic condition capable of causing widespread damage or loss of property and life. These hazards are geological and environmental conditions and involve long-term or short-term geological processes. Geohazards can be relatively small features, but they can also attain huge dimensions (e.g., submarine or surface landslide) and affect local and regional socio-economics to a large extent (e.g., tsunamis).
Sometimes the hazard is instigated by the careless location of developments or construction in which the conditions were not taken into account. Human activities, such as drilling through overpressured zones, could result in significant risk, and as such mitigation and prevention are paramount, through improved understanding of geohazards, their preconditions, causes and implications. In other cases, particularly in montane regions, natural processes can cause catalytic events of a complex nature, such as an avalanche hitting a lake and causing a debris flow, with consequences potentially hundreds of miles away, or creating a lahar by volcanism.
Marine geohazards in particular constitute a fast-growing sector of research as they involve seismic, tectonic, volcanic processes now occurring at higher frequency, and often resulting in coastal sub-marine avalanches or devastating tsunamis in some of the most densely populated areas of the world
Such impacts on vulnerable coastal populations, coastal infrastructures, offshore exploration platforms, obviously call for a higher level of preparedness and mitigation.
Speed of development
Sudden phenomena
Sudden phenomena include:
avalanches (snow or rock) and its runout
earthquakes and earthquake-triggered phenomena such as tsunamis
forest fires (espec. in Mediterranean areas) leading to deforestation
geomagnetic storms
gulls (chasms) associated with cambering of valley sides
ice jams (Eisstoß) on rivers or glacial lake outburst floods below a glacier
landslide (displacement of earth materials on a slope or hillside)
mudflows (avalanche-like muddy flow of soft/wet soil and sediment materials, narrow landslides)
pyroclastic flows
rockfalls, rock slides, (rock avalanche) and debris flows
torrents (flash floods, rapid floods or heavy current creeks with irregular course)
liquefaction (settlement of the ground in areas underlain by loose saturated sand/silt during an earthquake event)
volcanic eruptions, lahars and ash falls.
Slow phenomena
Gradual or slow phenomena include:
alluvial fans (e.g. at the exit of canyons or side valleys)
caldera development (volcanoes)
geyser deposits
ground settlement due to consolidation of compressible soils or due to collapseable soils (''see also compaction)
ground subsidence, sags and sinkholes
sand dune migration
shoreline and stream erosion
thermal springs
Evaluation and mitigation
Geologic hazards are typically evaluated by engineering geologists who are educated and trained in interpretation of landforms and earth process, earth-structure interaction, and in geologic hazard mitigation. The engineering geologist provides recommendations and designs to mitigate for geologic hazards. Trained hazard mitigation planners also assist local communities to identify strategies for mitigating the effects of such hazards and developing plans to implement these measures. Mitigation can include a variety of measures:
Geologic hazards may be avoided by relocation. Publicly available databases, via searchable platforms, can help people evaluate hazards in locations of interest.
Mapping geohazards using conventional or remote sensing techniques can also help identify suitable areas for urban development.
The stability of sloping earth can be improved by the construction of retaining walls, which may use techniques such as slurry walls, shear pins, tiebacks, soil nails or soil anchors. Larger projects may use gabions and other forms of earth buttress.
Shorelines and streams are protected against scour and erosion using revetments and riprap.
The soil or rock itself may be improved by means such as dynamic compaction, injection of grout or concrete, and mechanically stabilized earth.
Additional mitigation methods include deep foundations, tunnels, surface and subdrain systems, and other measures.
Planning measures include regulations prohibiting development near hazard-prone areas and adoption of building codes.
Earth observation of geohazards
In recent decades, Earth Observation (EO) has become a key tool in geohazards management, including preparedness, response, recovery, and mitigation. By leveraging remote sensing technologies, often supported by ground surveys, EO provides critical information to researchers, decision-makers, and planners. It has revolutionized our ability to map and monitor geohazards with precision and timeliness.
In paleohistory
Eleven distinct flood basalt episodes occurred in the past 250 million years, resulting in large volcanic provinces, creating lava plateaus and mountain ranges on Earth. Large igneous provinces have been connected to five mass extinction events. The timing of six out of eleven known provinces coincide with periods of global warming and marine anoxia/dysoxia. Thus, suggesting that volcanic CO2 emissions can force an important effect on the climate system.
Known hazards
2004 Indian Ocean earthquake and tsunami
2008 Sichuan earthquake
2011 Tōhoku earthquake and tsunami
The Barrier (located in Garibaldi Provincial Park)
Usoi Dam a natural landslide dam
| Physical sciences | Geology: General | Earth science |
782427 | https://en.wikipedia.org/wiki/Cone | Cone | A cone is a three-dimensional geometric shape that tapers smoothly from a flat base (frequently, though not necessarily, circular) to a point called the apex or vertex that is not contained in the base.
A cone is formed by a set of line segments, half-lines, or lines connecting a common point, the apex, to all of the points on a base.
In the case of line segments, the cone does not extend beyond the base, while in the case of half-lines, it extends infinitely far. In the case of lines, the cone extends infinitely far in both directions from the apex, in which case it is sometimes called a double cone. Each of the two halves of a double cone split at the apex is called a nappe.
Depending on the author, the base may be restricted to be a circle, any one-dimensional quadratic form in the plane, any closed one-dimensional figure, or any of the above plus all the enclosed points. If the enclosed points are included in the base, the cone is a solid object; otherwise it is an open surface, a two-dimensional object in three-dimensional space. In the case of a solid object, the boundary formed by these lines or partial lines is called the lateral surface; if the lateral surface is unbounded, it is a conical surface.
The axis of a cone is the straight line passing through the apex about which the cone has a circular symmetry.
In common usage in elementary geometry, cones are assumed to be right circular, i.e., with a circle base perpendicular to the axis. If the cone is right circular the intersection of a plane with the lateral surface is a conic section. In general, however, the base may be any shape and the apex may lie anywhere (though it is usually assumed that the base is bounded and therefore has finite area, and that the apex lies outside the plane of the base). Contrasted with right cones are oblique cones, in which the axis passes through the centre of the base non-perpendicularly.
Depending on the context, "cone" may also mean specifically a convex cone or a projective cone.
Cones can also be generalized to higher dimensions.
Further terminology
The perimeter of the base of a cone is called the "directrix", and each of the line segments between the directrix and apex is a "generatrix" or "generating line" of the lateral surface. (For the connection between this sense of the term "directrix" and the directrix of a conic section, see Dandelin spheres.)
The "base radius" of a circular cone is the radius of its base; often this is simply called the radius of the cone. The aperture of a right circular cone is the maximum angle between two generatrix lines; if the generatrix makes an angle θ to the axis, the aperture is 2θ. In optics, the angle θ is called the half-angle of the cone, to distinguish it from the aperture.
A cone with a region including its apex cut off by a plane is called a truncated cone; if the truncation plane is parallel to the cone's base, it is called a frustum. An elliptical cone is a cone with an elliptical base. A generalized cone is the surface created by the set of lines passing through a vertex and every point on a boundary (also see visual hull).
Measurements and equations
Volume
The volume of any conic solid is one third of the product of the area of the base and the height
In modern mathematics, this formula can easily be computed using calculus — it is, up to scaling, the integral Without using calculus, the formula can be proven by comparing the cone to a pyramid and applying Cavalieri's principle – specifically, comparing the cone to a (vertically scaled) right square pyramid, which forms one third of a cube. This formula cannot be proven without using such infinitesimal arguments – unlike the 2-dimensional formulae for polyhedral area, though similar to the area of the circle – and hence admitted less rigorous proofs before the advent of calculus, with the ancient Greeks using the method of exhaustion. This is essentially the content of Hilbert's third problem – more precisely, not all polyhedral pyramids are scissors congruent (can be cut apart into finite pieces and rearranged into the other), and thus volume cannot be computed purely by using a decomposition argument.
Center of mass
The center of mass of a conic solid of uniform density lies one-quarter of the way from the center of the base to the vertex, on the straight line joining the two.
Right circular cone
Volume
For a circular cone with radius r and height h, the base is a circle of area and so the formula for volume becomes
Slant height
The slant height of a right circular cone is the distance from any point on the circle of its base to the apex via a line segment along the surface of the cone. It is given by , where is the radius of the base and is the height. This can be proved by the Pythagorean theorem.
Surface area
The lateral surface area of a right circular cone is where is the radius of the circle at the bottom of the cone and is the slant height of the cone. The surface area of the bottom circle of a cone is the same as for any circle, . Thus, the total surface area of a right circular cone can be expressed as each of the following:
Radius and height
(the area of the base plus the area of the lateral surface; the term is the slant height)
where is the radius and is the height.
Radius and slant height
where is the radius and is the slant height.
Circumference and slant height
where is the circumference and is the slant height.
Apex angle and height
where is the apex angle and is the height.
Circular sector
The circular sector is obtained by unfolding the surface of one nappe of the cone:
radius R
arc length L
central angle φ in radians
Equation form
The surface of a cone can be parameterized as
where is the angle "around" the cone, and is the "height" along the cone.
A right solid circular cone with height and aperture , whose axis is the coordinate axis and whose apex is the origin, is described parametrically as
where range over , , and , respectively.
In implicit form, the same solid is defined by the inequalities
where
More generally, a right circular cone with vertex at the origin, axis parallel to the vector , and aperture , is given by the implicit vector equation where
where , and denotes the dot product.
Elliptic cone
In the Cartesian coordinate system, an elliptic cone is the locus of an equation of the form
It is an affine image of the right-circular unit cone with equation From the fact, that the affine image of a conic section is a conic section of the same type (ellipse, parabola,...), one gets:
Any plane section of an elliptic cone is a conic section.
Obviously, any right circular cone contains circles. This is also true, but less obvious, in the general case (see circular section).
The intersection of an elliptic cone with a concentric sphere is a spherical conic.
Projective geometry
In projective geometry, a cylinder is simply a cone whose apex is at infinity. Intuitively, if one keeps the base fixed and takes the limit as the apex goes to infinity, one obtains a cylinder, the angle of the side increasing as arctan, in the limit forming a right angle. This is useful in the definition of degenerate conics, which require considering the cylindrical conics.
According to G. B. Halsted, a cone is generated similarly to a Steiner conic only with a projectivity and axial pencils (not in perspective) rather than the projective ranges used for the Steiner conic:
"If two copunctual non-costraight axial pencils are projective but not perspective, the meets of correlated planes form a 'conic surface of the second order', or 'cone'."
Generalizations
The definition of a cone may be extended to higher dimensions; see convex cone. In this case, one says that a convex set C in the real vector space is a cone (with apex at the origin) if for every vector x in C and every nonnegative real number a, the vector ax is in C. In this context, the analogues of circular cones are not usually special; in fact one is often interested in polyhedral cones.
An even more general concept is the topological cone, which is defined in arbitrary topological spaces.
| Mathematics | Three-dimensional space | null |
22284121 | https://en.wikipedia.org/wiki/Variable%20%28computer%20science%29 | Variable (computer science) | In computer programming, a variable is an abstract storage location paired with an associated symbolic name, which contains some known or unknown quantity of data or object referred to as a value; or in simpler terms, a variable is a named container for a particular set of bits or type of data (like integer, float, string, etc...). A variable can eventually be associated with or identified by a memory address. The variable name is the usual way to reference the stored value, in addition to referring to the variable itself, depending on the context. This separation of name and content allows the name to be used independently of the exact information it represents. The identifier in computer source code can be bound to a value during run time, and the value of the variable may thus change during the course of program execution.
Variables in programming may not directly correspond to the concept of variables in mathematics. The latter is abstract, having no reference to a physical object such as storage location. The value of a computing variable is not necessarily part of an equation or formula as in mathematics. Variables in computer programming are frequently given long names to make them relatively descriptive of their use, whereas variables in mathematics often have terse, one- or two-character names for brevity in transcription and manipulation.
A variable's storage location may be referenced by several different identifiers, a situation known as aliasing. Assigning a value to the variable using one of the identifiers will change the value that can be accessed through the other identifiers.
Compilers have to replace variables' symbolic names with the actual locations of the data. While a variable's name, type, and location often remain fixed, the data stored in the location may be changed during program execution.
Actions on a variable
In imperative programming languages, values can generally be accessed or changed at any time. In pure functional and logic languages, variables are bound to expressions and keep a single value during their entire lifetime due to the requirements of referential transparency. In imperative languages, the same behavior is exhibited by (named) constants (symbolic constants), which are typically contrasted with (normal) variables.
Depending on the type system of a programming language, variables may only be able to store a specified data type (e.g. integer or string). Alternatively, a datatype may be associated only with the current value, allowing a single variable to store anything supported by the programming language. Variables are the containers for storing the values.
Variables and scope:
Automatic variables: Each local variable in a function comes into existence only when the function is called, and disappears when the function is exited. Such variables are known as automatic variables.
External variables: These are variables that are external to a function and can be accessed by name by any function. These variables remain in existence permanently; rather than appearing and disappearing as functions are called and exited, they retain their values even after the functions that set them have returned.
Identifiers referencing a variable
An identifier referencing a variable can be used to access the variable in order to read out the value, or alter the value, or edit other attributes of the variable, such as access permission, locks, semaphores, etc.
For instance, a variable might be referenced by the identifier "" and the variable can contain the number 1956. If the same variable is referenced by the identifier "" as well, and if using this identifier "", the value of the variable is altered to 2009, then reading the value using the identifier "" will yield a result of 2009 and not 1956.
If a variable is only referenced by a single identifier, that identifier can simply be called the name of the variable; otherwise, we can speak of it as one of the names of the variable. For instance, in the previous example the identifier "" is the name of the variable in question, and "" is another name of the same variable.
Scope and extent
The scope of a variable describes where in a program's text the variable may be used, while the extent (also called lifetime) of a variable describes when in a program's execution the variable has a (meaningful) value. The scope of a variable affects its extent. The scope of a variable is actually a property of the name of the variable, and the extent is a property of the storage location of the variable. These should not be confused with context (also called environment), which is a property of the program, and varies by point in the program's text or execution—see scope: an overview. Further, object lifetime may coincide with variable lifetime, but in many cases is not tied to it.
Scope is an important part of the name resolution of a variable. Most languages define a specific scope for each variable (as well as any other named entity), which may differ within a given program. The scope of a variable is the portion of the program's text for which the variable's name has meaning and for which the variable is said to be "visible". Entrance into that scope typically begins a variable's lifetime (as it comes into context) and exit from that scope typically ends its lifetime (as it goes out of context). For instance, a variable with "lexical scope" is meaningful only within a certain function/subroutine, or more finely within a block of expressions/statements (accordingly with function scope or block scope); this is static resolution, performable at parse-time or compile-time. Alternatively, a variable with dynamic scope is resolved at run-time, based on a global binding stack that depends on the specific control flow. Variables only accessible within a certain functions are termed "local variables". A "global variable", or one with indefinite scope, may be referred to anywhere in the program.
Extent, on the other hand, is a runtime (dynamic) aspect of a variable. Each binding of a variable to a value can have its own extent at runtime. The extent of the binding is the portion of the program's execution time during which the variable continues to refer to the same value or memory location. A running program may enter and leave a given extent many times, as in the case of a closure.
Unless the programming language features garbage collection, a variable whose extent permanently outlasts its scope can result in a memory leak, whereby the memory allocated for the variable can never be freed since the variable which would be used to reference it for deallocation purposes is no longer accessible. However, it can be permissible for a variable binding to extend beyond its scope, as occurs in Lisp closures and C static local variables; when execution passes back into the variable's scope, the variable may once again be used. A variable whose scope begins before its extent does is said to be uninitialized and often has an undefined, arbitrary value if accessed (see wild pointer), since it has yet to be explicitly given a particular value. A variable whose extent ends before its scope may become a dangling pointer and deemed uninitialized once more since its value has been destroyed. Variables described by the previous two cases may be said to be out of extent or unbound. In many languages, it is an error to try to use the value of a variable when it is out of extent. In other languages, doing so may yield unpredictable results. Such a variable may, however, be assigned a new value, which gives it a new extent.
For space efficiency, a memory space needed for a variable may be allocated only when the variable is first used and freed when it is no longer needed. A variable is only needed when it is in scope, thus beginning each variable's lifetime when it enters scope may give space to unused variables. To avoid wasting such space, compilers often warn programmers if a variable is declared but not used.
It is considered good programming practice to make the scope of variables as narrow as feasible so that different parts of a program do not accidentally interact with each other by modifying each other's variables. Doing so also prevents action at a distance. Common techniques for doing so are to have different sections of a program use different name spaces, or to make individual variables "private" through either dynamic variable scoping or lexical variable scoping.
Many programming languages employ a reserved value (often named null or nil) to indicate an invalid or uninitialized variable.
Typing
In statically typed languages such as C, C++, Java or C#, a variable also has a type, meaning that only certain kinds of values can be stored in it. For example, a variable of type "integer" is prohibited from storing text values.
In dynamically typed languages such as Python, a variable's type is inferred by its value, and can change according to its value. In Common Lisp, both situations exist simultaneously: A variable is given a type (if undeclared, it is assumed to be , the universal supertype) which exists at compile time. Values also have types, which can be checked and queried at runtime.
Typing of variables also allows polymorphisms to be resolved at compile time. However, this is different from the polymorphism used in object-oriented function calls (referred to as virtual functions in C++) which resolves the call based on the value type as opposed to the supertypes the variable is allowed to have.
Variables often store simple data, like integers and literal strings, but some programming languages allow a variable to store values of other datatypes as well. Such languages may also enable functions to be parametric polymorphic. These functions operate like variables to represent data of multiple types. For example, a function named may determine the length of a list. Such a function may be parametric polymorphic by including a type variable in its type signature, since the number of elements in the list is independent of the elements' types.
Parameters
The formal parameters (or formal arguments) of functions are also referred to as variables. For instance, in this Python code segment,
>>> def addtwo(x):
... return x + 2
...
>>> addtwo(5)
7
the variable named is a parameter because it is given a value when the function is called. The integer 5 is the argument which gives its value. In most languages, function parameters have local scope. This specific variable named can only be referred to within the function (though of course other functions can also have variables called ).
Memory allocation
The specifics of variable allocation and the representation of their values vary widely, both among programming languages and among implementations of a given language. Many language implementations allocate space for local variables, whose extent lasts for a single function call on the call stack, and whose memory is automatically reclaimed when the function returns. More generally, in name binding, the name of a variable is bound to the address of some particular block (contiguous sequence) of bytes in memory, and operations on the variable manipulate that block. Referencing is more common for variables whose values have large or unknown sizes when the code is compiled. Such variables reference the location of the value instead of storing the value itself, which is allocated from a pool of memory called the heap.
Bound variables have values. A value, however, is an abstraction, an idea; in implementation, a value is represented by some data object, which is stored somewhere in computer memory. The program, or the runtime environment, must set aside memory for each data object and, since memory is finite, ensure that this memory is yielded for reuse when the object is no longer needed to represent some variable's value.
Objects allocated from the heap must be reclaimed—especially when the objects are no longer needed. In a garbage-collected language (such as C#, Java, Python, Golang and Lisp), the runtime environment automatically reclaims objects when extant variables can no longer refer to them. In non-garbage-collected languages, such as C, the program (and the programmer) must explicitly allocate memory, and then later free it, to reclaim its memory. Failure to do so leads to memory leaks, in which the heap is depleted as the program runs, risks eventual failure from exhausting available memory.
When a variable refers to a data structure created dynamically, some of its components may be only indirectly accessed through the variable. In such circumstances, garbage collectors (or analogous program features in languages that lack garbage collectors) must deal with a case where only a portion of the memory reachable from the variable needs to be reclaimed.
Naming conventions
Unlike their mathematical counterparts, programming variables and constants commonly take multiple-character names, e.g. or . Single-character names are most commonly used only for auxiliary variables; for instance, , , for array index variables.
Some naming conventions are enforced at the language level as part of the language syntax which involves the format of valid identifiers. In almost all languages, variable names cannot start with a digit (0–9) and cannot contain whitespace characters. Whether or not punctuation marks are permitted in variable names varies from language to language; many languages only permit the underscore ("_") in variable names and forbid all other punctuation. In some programming languages, sigils (symbols or punctuation) are affixed to variable identifiers to indicate the variable's datatype or scope.
Case-sensitivity of variable names also varies between languages and some languages require the use of a certain case in naming certain entities; Most modern languages are case-sensitive; some older languages are not. Some languages reserve certain forms of variable names for their own internal use; in many languages, names beginning with two underscores ("__") often fall under this category.
However, beyond the basic restrictions imposed by a language, the naming of variables is largely a matter of style. At the machine code level, variable names are not used, so the exact names chosen do not matter to the computer. Thus names of variables identify them, for the rest they are just a tool for programmers to make programs easier to write and understand. Using poorly chosen variable names can make code more difficult to review than non-descriptive names, so names that are clear are often encouraged.
Programmers often create and adhere to code style guidelines that offer guidance on naming variables or impose a precise naming scheme. Shorter names are faster to type but are less descriptive; longer names often make programs easier to read and the purpose of variables easier to understand. However, extreme verbosity in variable names can also lead to less comprehensible code.
Variable types (based on lifetime)
We can classify variables based on their lifetime. The different types of variables are static, stack-dynamic, explicit heap-dynamic, and implicit heap-dynamic. A static variable is also known as global variable, it is bound to a memory cell before execution begins and remains to the same memory cell until termination. A typical example is the static variables in C and C++. A Stack-dynamic variable is known as local variable, which is bound when the declaration statement is executed, and it is deallocated when the procedure returns. The main examples are local variables in C subprograms and Java methods. Explicit Heap-Dynamic variables are nameless (abstract) memory cells that are allocated and deallocated by explicit run-time instructions specified by the programmer. The main examples are dynamic objects in C++ (via new and delete) and all objects in Java. Implicit Heap-Dynamic variables are bound to heap storage only when they are assigned values. Allocation and release occur when values are reassigned to variables. As a result, Implicit heap-dynamic variables have the highest degree of flexibility. The main examples are some variables in JavaScript, PHP and all variables in APL.
| Technology | Software development: General | null |
22284647 | https://en.wikipedia.org/wiki/Real-valued%20function | Real-valued function | In mathematics, a real-valued function is a function whose values are real numbers. In other words, it is a function that assigns a real number to each member of its domain.
Real-valued functions of a real variable (commonly called real functions) and real-valued functions of several real variables are the main object of study of calculus and, more generally, real analysis. In particular, many function spaces consist of real-valued functions.
Algebraic structure
Let be the set of all functions from a set to real numbers . Because is a field, may be turned into a vector space and a commutative algebra over the reals with the following operations:
– vector addition
– additive identity
– scalar multiplication
– pointwise multiplication
These operations extend to partial functions from to with the restriction that the partial functions and are defined only if the domains of and have a nonempty intersection; in this case, their domain is the intersection of the domains of and .
Also, since is an ordered set, there is a partial order
on which makes a partially ordered ring.
Measurable
The σ-algebra of Borel sets is an important structure on real numbers. If has its σ-algebra and a function is such that the preimage of any Borel set belongs to that σ-algebra, then is said to be measurable. Measurable functions also form a vector space and an algebra as explained above in .
Moreover, a set (family) of real-valued functions on can actually define a σ-algebra on generated by all preimages of all Borel sets (or of intervals only, it is not important). This is the way how σ-algebras arise in (Kolmogorov's) probability theory, where real-valued functions on the sample space are real-valued random variables.
Continuous
Real numbers form a topological space and a complete metric space. Continuous real-valued functions (which implies that is a topological space) are important in theories of topological spaces and of metric spaces. The extreme value theorem states that for any real continuous function on a compact space its global maximum and minimum exist.
The concept of metric space itself is defined with a real-valued function of two variables, the metric, which is continuous. The space of continuous functions on a compact Hausdorff space has a particular importance. Convergent sequences also can be considered as real-valued continuous functions on a special topological space.
Continuous functions also form a vector space and an algebra as explained above in , and are a subclass of measurable functions because any topological space has the σ-algebra generated by open (or closed) sets.
Smooth
Real numbers are used as the codomain to define smooth functions. A domain of a real smooth function can be the real coordinate space (which yields a real multivariable function), a topological vector space, an open subset of them, or a smooth manifold.
Spaces of smooth functions also are vector spaces and algebras as explained above in and are subspaces of the space of continuous functions.
Appearances in measure theory
A measure on a set is a non-negative real-valued functional on a σ-algebra of subsets. Lp spaces on sets with a measure are defined from aforementioned real-valued measurable functions, although they are actually quotient spaces. More precisely, whereas a function satisfying an appropriate summability condition defines an element of Lp space, in the opposite direction for any and which is not an atom, the value is undefined. Though, real-valued Lp spaces still have some of the structure described above in . Each of Lp spaces is a vector space and have a partial order, and there exists a pointwise multiplication of "functions" which changes , namely
For example, pointwise product of two L2 functions belongs to L1.
Other appearances
Other contexts where real-valued functions and their special properties are used include monotonic functions (on ordered sets), convex functions (on vector and affine spaces), harmonic and subharmonic functions (on Riemannian manifolds), analytic functions (usually of one or more real variables), algebraic functions (on real algebraic varieties), and polynomials (of one or more real variables).
| Mathematics | Functions: General | null |
18575446 | https://en.wikipedia.org/wiki/Civil%20aviation | Civil aviation | Civil aviation is one of two major categories of flying, representing all non-military and non-state aviation, which can be both private and commercial. Most countries in the world are members of the International Civil Aviation Organization and work together to establish common Standards and Recommended Practices for civil aviation through that agency.
Civil aviation includes three major categories:
Commercial air transport, including scheduled and non-scheduled passenger and cargo flights
Aerial work, in which an aircraft is used for specialized services such as agriculture, photography, surveying, search and rescue, etc.
General aviation (GA), including all other civil flights, private or commercial
Although scheduled air transport is the larger operation in terms of passenger numbers, GA is larger in the number of flights (and flight hours, in the U.S.) In the U.S., GA carries 166 million passengers each year, more than any individual airline, though less than all the airlines combined. Since 2004, the U.S. airlines combined have carried over 600 million passengers each year, and in 2014, they carried a combined 662,819,232 passengers.
Some countries also make a regulatory distinction based on whether aircraft are flown for hire, like:
Commercial aviation includes most or all flying done for hire, particularly scheduled service on airlines; and
Private aviation includes pilots flying for their own purposes (recreation, business meetings, etc.) without receiving any kind of remuneration.
All scheduled air transport is commercial, but general aviation can be either commercial or private. Normally, the pilot, aircraft, and operator must all be authorized to perform commercial operations through separate commercial licensing, registration, and operation certificates.
Non-civil aviation is referred to as state aviation. This includes military aviation, state VIP transports, and police/customs aircraft.
History
Postwar aviation
After World War II, commercial aviation grew rapidly, using mostly ex-military pilots to transport people and cargo. Factories that had produced bombers were quickly adapted to the production of passenger aircraft like the Douglas DC-4. This growth was accelerated by the establishment of military airports throughout the world, either for combat use or training. These could easily be turned to civil aviation use. The first commercial jet airliner to fly was the British de Havilland DH.106 Comet. By 1952, the British state airline British Overseas Airways Corporation had introduced the Comet into scheduled service. While it was a technical achievement, the airplane suffered a series of highly public failures, as the shape of the windows led to cracks due to metal fatigue. By the time the problems were overcome, other jet airliner designs such as the Boeing 707 had already entered service.
Civil aviation authorities
The Chicago Convention on International Civil Aviation was originally established in 1944; it states that signatories should collectively work to harmonize and standardize the use of airspace for safety, efficiency and regularity of air transport. Each signatory country, of which there are at least 193, has a civil aviation authority (such as the Federal Aviation Administration in the United States) to oversee the following areas of civil aviation:
Personnel licensing — regulating the basic training and issuance of licenses and certificates.
Flight operations — carrying out safety oversight of commercial operators.
Airworthiness — issuing certificates of registration and certificates of airworthiness to civil aircraft, and overseeing the safety of aircraft maintenance organizations.
Aerodromes — designing and constructing aerodrome facilities.
Air traffic services — managing the traffic inside of a country's airspace.
Statistics
The World Bank lists monotonously growing numbers for the number of passengers transported per year worldwide with a preliminary all-time high in 2015 of 3.44 billion passengers. Likewise, the number of registered carrier departures worldwide has reached a peak in 2015 with almost 33 million takeoffs. In the U.S. alone, the passenger miles "computed by summing the products of the aircraft-miles flown on each inter airport segment multiplied by the number of passengers carried on that segment" have reached in 2014 (as compared to highway car traffic with ). The global seasonally adjusted revenue passenger kilometers per month peaked at more than (~6.6 trillion per year, corresponding to roughly 2000 km per passenger) in January 2016, a 7% rise over one year. The passenger numbers are distinctively more volatile than general economic indicators. Global political, economic or health crises have an amplifying effect.
| Technology | Concepts of aviation | null |
23785293 | https://en.wikipedia.org/wiki/Mensa%20%28constellation%29 | Mensa (constellation) | Mensa is a constellation in the Southern Celestial Hemisphere near the south celestial pole, one of fourteen constellations drawn up in the 18th century by French astronomer Nicolas-Louis de Lacaille. Its name is Latin for table, though it originally commemorated Table Mountain and was known as "Mons Mensae". One of the eighty-eight constellations designated by the International Astronomical Union (IAU), it covers a keystone-shaped wedge of sky 153.5 square degrees in area. Other than the south polar constellation of Octans, it is the most southerly of constellations and is observable only south of the 5th parallel of the Northern Hemisphere.
One of the faintest constellations in the night sky, Mensa contains no apparently bright stars—the brightest, Alpha Mensae, is barely visible in suburban skies. Part of the Large Magellanic Cloud, several star clusters and a quasar lie in the area covered by the constellation, and at least three of its star systems have been found to have exoplanets.
History
Originally named "Montagne de la Table" or "Mons Mensae", Mensa was created by Nicolas-Louis de Lacaille out of dim Southern Hemisphere stars in honor of Table Mountain, a South African mountain overlooking Cape Town, near the location of Lacaille's observatory. He recalled that the Magellanic Clouds were sometimes known as Cape clouds, and that Table Mountain was often covered in clouds when a southeasterly stormy wind blew. Hence he made a "table" in the sky under the clouds. Lacaille had observed and catalogued 10,000 southern stars during a two-year stay at the Cape of Good Hope. He devised 14 new constellations in uncharted regions of the Southern Celestial Hemisphere not visible from Europe. Mensa was the only constellation that did not honor an instrument symbolic of the Age of Enlightenment. Sir John Herschel proposed shrinking the name to one word in 1844, noting that Lacaille himself had abbreviated some of his constellations thus.
Although the stars of Mensa do not feature in any ancient mythology, the mountain it is named after has a rich mythology. Called "Tafelberg" in Dutch and German, it has two neighboring mountains called "Devil's Peak" and "Lion's Head". Table Mountain features in the mythology of the Cape of Good Hope, notorious for its storms. Explorer Bartolomeu Dias saw the mountain as a mythical anvil for storms.
Characteristics
Mensa is bordered by Dorado to the north, Hydrus to the northwest and west, Octans to the south, Chamaeleon to the east and Volans to the northeast. Covering 153.5 square degrees and 0.372% of the night sky, it ranks 75th of the 88 constellations in size. The three-letter abbreviation for the constellation, as adopted by the IAU in 1922, is "Men". The official constellation boundaries, as set by Belgian astronomer Eugène Delporte in 1930, are defined by a polygon of eight segments. In the equatorial coordinate system, the right ascension coordinates of these borders lie between and , while the declination coordinates are between −69.75° and −85.26°. The whole constellation is visible to observers south of latitude 5°N.
Features
Stars
Bright stars
Lacaille gave eleven stars in the constellation Bayer designations, using the Greek alphabet to label them Alpha through to Lambda Mensae (excluding Kappa). Gould later added Kappa, Mu, Nu, Xi and Pi Mensae. Stars as dim as these were not generally given designations; however, Gould felt their closeness to the South Celestial Pole warranted their naming. Alpha Mensae is the brightest star with a barely visible apparent magnitude of 5.09, making it the only constellation with no star above magnitude 5.0. Overall, there are 22 stars within the constellation's borders brighter than or equal to apparent magnitude 6.5.
Alpha Mensae is a solar-type star (class G7V) 33.32 ± 0.02 light-years from Earth. It came to within 11 light-years from Earth around 250,000 years ago and would have been considerably brighter back then—nearly of second magnitude. An infrared excess has been detected around this star, indicating the presence of a circumstellar disk at a radius of over 147 astronomical units (AU). The estimated temperature of this dust is below 22 K. However, data from Herschel Space Observatory failed to confirm this excess, leaving the finding in doubt. No planetary companions have yet been discovered around it. It has a red dwarf companion star at an angular separation of 3.05 arcseconds; equivalent to a projected separation of roughly 30 AU.
Gamma Mensae is the second-brightest star in the constellation, at magnitude 5.19. Located 104.9 ± 0.5 light-years from Earth, it is an ageing (10.6 billion year-old) star around 1.04 times as massive as the Sun. It has swollen to around 5 times the solar radius, becoming an orange giant of spectral type K2III.
Beta Mensae is slightly fainter at magnitude 5.31. Located 660 ± 10 light-years from Earth, it is a yellow giant of spectral type G8III, around 3.6 times as massive and 513 times as luminous as the Sun. It is 270 million years old, and lies in front of the Large Magellanic Cloud.
Zeta and Eta Mensae have infrared excesses suggesting they too have circumstellar disks of dust. Zeta Mensae is an ageing white giant of spectral type A5 III around 394 ± 4 light-years from Earth, and Eta Mensae is an orange giant of spectral type K4 III, lying 650 ± 10 light-years away from Earth.
Pi Mensae is a solar-type (G1) star 59.62 ± 0.07 light-years distant. In 2001, a substellar companion was discovered in an eccentric orbit. Incorporating more accurate Hipparcos data yields a mass range for the companion to be anywhere from 10.27 to 29.9 times that of Jupiter. This confirms its substellar nature with the upper limit of mass putting it in the brown dwarf range. The discovery of a second substellar companion—a super-Earth—was announced on 16 September 2018. It takes 6.27 days to complete its orbit and is the first exoplanet detected by the Transiting Exoplanet Survey Satellite (TESS) submitted for publication.
Planet-hosting stars
HD 38283 (Bubup) is a Sun-like star of spectral type F9.5V of magnitude 6.7, located 124.3 ± 0.1 light-years distant. In 2011, a gas giant with an Earth-like orbital period of 363 days and a minimum mass a third that of Jupiter was discovered by the radial velocity method.
HD 39194 is an orange dwarf of spectral type K0V and magnitude 8.08, located 86.21 ± 0.09 light-years distant. Three planets in close orbit were discovered by the High Accuracy Radial Velocity Planet Searcher (HARPS) in 2011. The three take 5.6, 14 and 34 days to complete an orbit around their star, and have minimum masses 3.72, 5.94 and 5.14 times that of the Earth respectively.
Variable stars
TZ Mensae is an eclipsing binary that varies between magnitude 6.2 and 6.9 every 8.57 days. It is composed of two white main sequence stars in close orbit. One of these is of spectral type A0V, has a radius twice as that of the Sun and is 2.5 times as massive. The other, an A8V spectral type, has a radius 1.4 times that of the Sun and is 1.5 times as massive.
UX Mensae is another eclipsing binary system composed of two young stars around 1.2 times as massive as the Sun and 2.2 ± 0.5 billion years of age, orbiting each other every 4.19 days. The system is 338.2 ± 0.9 light-years distant.
TY Mensae is another eclipsing binary system classified as a W Ursae Majoris variable; the two components are so close that they share a common envelope of stellar material. The larger star has been calculated to be 1.86 times as massive, have 1.85 times the diameter and is 13.6 times as luminous, while the smaller is 0.4 times as massive, 0.84 times the diameter, and 1.7 times as luminous as the Sun. Their surface temperatures have been calculated at 8164 and 7183 K respectively.
YY Mensae is an orange giant of spectral type K1III around 2.2 times as massive as the Sun, with 12.7 times its diameter and 70 times its luminosity. A rapidly rotating star with a period of 9.5 days, it is a strong emitter of X-rays and belongs to a class of star known as FK Comae Berenices variables. These stars are thought to have formed with the merger of two stars in a contact binary system. With an apparent magnitude of 8.05, it is 707 ± 6 light-years distant.
AH Mensae is a cataclysmic variable star system composed of a white dwarf and a red dwarf that orbit each other every 2 hours 57 minutes. The stars are close enough that the white dwarf strips material off the red dwarf, creating an accretion disc that periodically ignites with a resulting brightening of the system.
TU Mensae is another cataclysmic variable composed of a red dwarf and white dwarf. The orbital period of 2 hours 49 minutes is one of the longest for cataclysmic variable systems exhibiting brighter outbursts, known as superhumps. The normal outbursts result in an increase in brightness lasting around a day every 37 days, while the superhumps last 5–20 days and take place every 194 days.
AO Mensae is a faint star of magnitude 9.8. An orange dwarf that has 80% the size and mass of the Sun, it is also a BY Draconis variable. These are a class of stars with starspots prominent enough that the star changes brightness as it rotates. It is a member of the Beta Pictoris moving group, a loose association of young stars moving across the galaxy.
Other stars
WISE 0535−7500 is a binary system composed of two sub-brown dwarfs of spectral class cooler than Y1 located 47 ± 3 light-years away. Unable to be separated by observations to date, they are presumed to be of similar mass—8 to 20 times that of Jupiter—and are less than one AU apart.
Deep-sky objects
The Large Magellanic Cloud lies partially within Mensa's boundaries, although most of it lies in neighbouring Dorado. It is a satellite galaxy of the Milky Way, located at a distance of 163,000 light-years. Among its stars within Mensa are W Mensae, an unusual yellow-white supergiant that belongs to a rare class of star known as a R Coronae Borealis variable, HD 268835, a blue hypergiant that is girded by a vast circumstellar disk of dust, and R71, a luminous blue variable star that brightened in 2012 to over a million times as luminous as the Sun. Also within the galaxy is NGC 1987, a globular cluster estimated to be around 600 million years old that has a significant number of red ageing stars, and NGC 1848, a 27 million year old open cluster. Mensa contains several described open clusters, most of which can be only be clearly observed from large telescopes.
PKS 0637-752 is a distant quasar with a calculated redshift of z = 0.651. It was chosen as the first target of the then newly-operational Chandra X-Ray Observatory in 1999. The resulting images revealed a gas jet approximately 330,000 light-years long. It is visible at radio, optical and x-ray wavelengths.
| Physical sciences | Other | Astronomy |
10575816 | https://en.wikipedia.org/wiki/Tropical%20cyclone%20forecasting | Tropical cyclone forecasting | Tropical cyclone forecasting is the science of forecasting where a tropical cyclone's center, and its effects, are expected to be at some point in the future. There are several elements to tropical cyclone forecasting: track forecasting, intensity forecasting, rainfall forecasting, storm surge, tornado, and seasonal forecasting. While skill is increasing in regard to track forecasting, intensity forecasting skill remains unchanged over the past several years. Seasonal forecasting began in the 1980s in the Atlantic basin and has spread into other basins in the years since.
History
Short term
The methods through which tropical cyclones are forecast have changed with the passage of time. The first known forecasts in the Western Hemisphere were made by Lt. Col. William Reed of the Corps of Royal Engineers at Barbados in 1847. Reed mostly utilized barometric pressure measurements as the basis of his forecasts. Benito Vines introduced a forecast and warning system based on cloud cover changes in Havana during the 1870s. Before the early 1900s, though, most forecasts were done by direct observations at weather stations, which were then relayed to forecast centres via telegraph. It wasn't until the advent of radio in the early twentieth century that observations from ships at sea were available to forecasters. The 1930s saw the usage of radiosondes in tropical cyclone forecasting. The next decade saw the advent of aircraft-based reconnaissance by the military, starting with the first dedicated flight into a hurricane in 1943, and the establishment of the Hurricane Hunters in 1944. In the 1950s, coastal weather radars began to be used in the United States, and research reconnaissance flights by the precursor of the Hurricane Research Division began in 1954.
The launch of the first weather satellite, TIROS-I, in 1960, introduced new forecasting techniques that remain important to tropical cyclone forecasting to the present. In the 1970s, buoys were introduced to improve the resolution of surface measurements, which until that point, were not available at all oversea surfaces.
Long term
In the late 1970s, William Gray noticed a trend of low hurricane activity in the North Atlantic basin during El Niño years. He was the first researcher to make a connection between such events and positive results led him to pursue further research. He found numerous factors across the globe influence tropical cyclone activity, such as connecting wet periods over the African Sahel to an increase in major hurricane landfalls along the United States East Coast. However, his findings also showed inconsistencies when only looking at a single factor as a primary influence.
Utilizing his findings, Gray developed an objective, statistical forecast for seasonal hurricane activity; he predicted only the number of tropical storms, hurricanes, and major hurricanes, foregoing specifics on tracks and potential landfalls due to the aforementioned inconsistencies. Gray issued his first seasonal forecast ahead of the 1984 season, which used the statistical relationships between tropical cyclone activity, the El Niño–Southern Oscillation (ENSO), Quasi-biennial oscillation (QBO), and Caribbean basin sea-level pressures. The endeavour proved modestly successful. He subsequently issued forecasts ahead of the start of the Atlantic hurricane season in May and before the peak of the season in August. Students and colleagues joined his forecast team in the following years, including Christopher Landsea, Paul W. Mielke Jr., and Kenneth J. Berry.
Track
The large-scale synoptic flow determines 70 to 90 percent of a tropical cyclone's motion. The deep-layer mean flow is the best tool in determining track direction and speed. If storms are significantly sheared, use of a lower-level wind is a better predictor. Knowledge of the beta effect can be used to steer a tropical cyclone, since it leads to a more northwest heading for tropical cyclones in the Northern Hemisphere. It is also best to smooth out short term wobbles of the storm centre to determine a more accurate trajectory.
Because of the forces that affect tropical cyclone tracks, accurate track predictions depend on determining the position and strength of high- and low-pressure areas and predicting how those areas will change during the life of a tropical system. Combining forecast models with increased understanding of the forces that act on tropical cyclones, and a wealth of data from Earth-orbiting satellites and other sensors, scientists have increased the accuracy of track forecasts over recent decades. An accurate track forecast is important, because if the track forecast is incorrect, forecasts for intensity, rainfall, storm surge, and tornado threat will also be incorrect.
1-2-3 rule
The 1-2-3 rule (mariner's 1-2-3 rule or danger area) is a guideline commonly taught to mariners for severe storm (specifically hurricane and tropical storm) tracking and prediction. The 1-2-3 rule has two parts, the 34-Knot Rule which is the danger area to be avoided. The 1-2-3 rule itself refers to the rounded long-term NHC/TPC forecast errors of 100-200-300 nautical miles at 24-48-72 hours, respectively. These numbers were close to the 10-year average for the 1982–1991-time frame. However, these errors have decreased to near 50-100-150 as NHC forecasters become more accurate. The "danger area" to be avoided is constructed by expanding the forecast path by a radius equal to the respective hundreds of miles plus the forecast 34-Knot wind field radii.
Intensity
Forecasters say they are less skilful at predicting the intensity of tropical cyclones than cyclone track. Available computing power limits forecasters' ability to accurately model a large number of complex factors, such as exact topology and atmospheric conditions, though with increased experience and understanding, even models with the same resolution can be tuned to more accurately reflect real-world behaviour. Another weakness is lack of frequent wind speed measurements in the eye of the storm. The Cyclone Global Navigation Satellite System, launched by NASA in 2016, is expected to provide much more data compared to sporadic measurements by weather buoys and hurricane-penetrating aircraft.
An accurate track forecast is essential to creating accurate intensity forecasts, particularly in an area with large islands such as the western north Pacific and the Caribbean Sea, as proximity to land is an inhibiting factor to developing tropical cyclones. A strong hurricane/typhoon/cyclone can weaken if an outer eye wall forms (typically around from the centre of the storm), choking off the convection within the inner eye wall. Such weakening is called an eyewall replacement cycle, and is usually temporary.
Maximum potential intensity
Dr. Kerry Emanuel created a mathematical model around 1988, called the maximum potential intensity or MPI, to compute the upper limit of tropical cyclone intensity based on sea surface temperature and atmospheric profiles from the latest global model runs. Maps created from this equation show values of the maximum achievable intensity due to the thermodynamics of the atmosphere at the time of the last model run (either 0000 or 1200 UTC). However, MPI does not take vertical wind shear into account. MPI is computed using the following formula:
Where is the maximum potential velocity in meters per second; is the sea surface temperature underneath the center of the tropical cyclone, is a reference temperature (30 °C) and , and are curve-fit constants. When , , and , the graph generated by this function corresponds to the 99th percentile of empirical tropical cyclone intensity data.
Rainfall
Tropical cyclone rainfall forecasting is important, since between 1970 and 2004, inland flooding from tropical cyclones caused most of the fatalities from tropical cyclones in the United States. While flooding is common to tropical cyclones near a landmass, there are a few factors which lead to excessive rainfall from tropical cyclones. Slow motion, as was seen during Hurricane Danny and Hurricane Wilma, can lead to high amounts. The presence of topography near the coast, as is the case across much of Mexico, Haiti, the Dominican Republic, much of Central America, Madagascar, Réunion, China, and Japan acts to magnify amounts due to upslope flow into the mountains. Strong upper level forcing from a trough moving through the Westerlies, as was the case during Hurricane Floyd, can lead to excessive amounts even from systems moving at an average forward motion. A combination of two of these factors could be especially crippling, as was seen during Hurricane Mitch in Central America. Therefore, an accurate track forecast is essential in order to produce an accurate tropical cyclone rainfall forecast.
However, as a result of global warming, the heat that has built up on the ocean's surface has allowed storms and hurricanes to capture more water vapour and, given the increased temperatures in the atmosphere also, retain the moisture for a longer capacity. This results in incredible amounts of rainfall upon striking land which can often be the most damaging aspect of a hurricane.
Operational methods
Historically, tropical cyclone tracking charts were used to include the past track and prepare future forecasts at Regional Specialized Meteorological Centers and Tropical Cyclone Warning Centers. The need for a more modernized method for forecasting tropical cyclones had become apparent to operational weather forecasters by the mid-1980s. At that time the United States Department of Defense was using paper maps, acetate, grease pencils, and disparate computer programs to forecast tropical cyclones. The Automated Tropical Cyclone Forecasting System (ATCF) software was developed by the Naval Research Laboratory for the Joint Typhoon Warning Center (JTWC) beginning in 1986, and used since 1988. During 1990 the system was adapted by the National Hurricane Center (NHC) for use at the NHC, National Centers for Environmental Prediction and the Central Pacific Hurricane Center. This provided the NHC with a multitasking software environment which allowed them to improve efficiency and cut the time required to make a forecast by 25% or 1 hour. ATCF was originally developed for use within DOS, before later being adapted to Unix and Linux.
Storm surge
The main storm surge forecast model in the Atlantic basin is SLOSH, which stands for Sea, Lake, Overland, Surge from Hurricanes. It uses the size of a storm, its intensity, its forward motion, and the topography of the coastal plain to estimate the depth of a storm surge at any individual grid point across the United States. An accurate forecast track is required in order to produce accurate storm surge forecasts. However, if the landfall point is uncertain, a maximum envelope of water (MEOW) map can be generated based on the direction of approach. If the forecast track itself is also uncertain, a maximum of maximums (MoM) map can be generated which will show the worst possible scenario for a hurricane of a specific strength.
Tornado
The location of most tropical cyclone-related tornadoes is their northeast quadrant in the Northern Hemisphere and southeast quadrant in the Southern Hemisphere. Like most of the other forecasts for tropical cyclone effects, an accurate track forecast is required in order to produce an accurate tornado threat forecast.
Seasonal forecast
By looking at annual variations in various climate parameters, forecasters can make predictions about the overall number and intensity of tropical cyclones that will occur in a given season. For example, when constructing its seasonal outlooks, the Climate Prediction Center in the United States considers the effects of the El Niño-Southern Oscillation, 25–40 year tropical cycle, wind shear over the oceans, and ocean surface temperature.
| Physical sciences | Storms | Earth science |
2585991 | https://en.wikipedia.org/wiki/Ehrenfest%20theorem | Ehrenfest theorem | The Ehrenfest theorem, named after Austrian theoretical physicist Paul Ehrenfest, relates the time derivative of the expectation values of the position and momentum operators x and p to the expectation value of the force on a massive particle moving in a scalar potential ,
The Ehrenfest theorem is a special case of a more general relation between the expectation of any quantum mechanical operator and the expectation of the commutator of that operator with the Hamiltonian of the system
where is some quantum mechanical operator and is its expectation value.
It is most apparent in the Heisenberg picture of quantum mechanics, where it amounts to just the expectation value of the Heisenberg equation of motion. It provides mathematical support to the correspondence principle.
The reason is that Ehrenfest's theorem is closely related to Liouville's theorem of Hamiltonian mechanics, which involves the Poisson bracket instead of a commutator. Dirac's rule of thumb suggests that statements in quantum mechanics which contain a commutator correspond to statements in classical mechanics where the commutator is supplanted by a Poisson bracket multiplied by . This makes the operator expectation values obey corresponding classical equations of motion, provided the Hamiltonian is at most quadratic in the coordinates and momenta. Otherwise, the evolution equations still may hold approximately, provided fluctuations are small.
Relation to classical physics
Although, at first glance, it might appear that the Ehrenfest theorem is saying that the quantum mechanical expectation values obey Newton’s classical equations of motion, this is not actually the case. If the pair were to satisfy Newton's second law, the right-hand side of the second equation would have to be
which is typically not the same as
If for example, the potential is cubic, (i.e. proportional to ), then is quadratic (proportional to ). This means, in the case of Newton's second law, the right side would be in the form of , while in the Ehrenfest theorem it is in the form of . The difference between these two quantities is the square of the uncertainty in and is therefore nonzero.
An exception occurs in case when the classical equations of motion are linear, that is, when is quadratic and is linear. In that special case, and do agree. Thus, for the case of a quantum harmonic oscillator, the expected position and expected momentum do exactly follow the classical trajectories.
For general systems, if the wave function is highly concentrated around a point , then and will be almost the same, since both will be approximately equal to . In that case, the expected position and expected momentum will approximately follow the classical trajectories, at least for as long as the wave function remains localized in position.
Derivation in the Schrödinger picture
Suppose some system is presently in a quantum state . If we want to know the instantaneous time derivative of the expectation value of , that is, by definition
where we are integrating over all of space. If we apply the Schrödinger equation, we find that
By taking the complex conjugate we find
Note , because the Hamiltonian is Hermitian. Placing this into the above equation we have
Often (but not always) the operator is time-independent so that its derivative is zero and we can ignore the last term.
Derivation in the Heisenberg picture
In the Heisenberg picture, the derivation is straightforward. The Heisenberg picture moves the time dependence of the system to operators instead of state vectors. Starting with the Heisenberg equation of motion,
Ehrenfest's theorem follows simply upon projecting the Heisenberg equation onto from the right and from the left, or taking the expectation value, so
One may pull the out of the first term, since the state vectors are no longer time dependent in the Heisenberg Picture. Therefore,
General example
For the very general example of a massive particle moving in a potential, the Hamiltonian is simply
where is the position of the particle.
Suppose we wanted to know the instantaneous change in the expectation of the momentum . Using Ehrenfest's theorem, we have
since the operator commutes with itself and has no time dependence. By expanding the right-hand-side, replacing by , we get
After applying the product rule on the second term, we have
As explained in the introduction, this result does not say that the pair satisfies Newton's second law, because the right-hand side of the formula is rather than . Nevertheless, as explained in the introduction, for states that are highly localized in space, the expected position and momentum will approximately follow classical trajectories, which may be understood as an instance of the correspondence principle.
Similarly, we can obtain the instantaneous change in the position expectation value.
This result is actually in exact accord with the classical equation.
Derivation of the Schrödinger equation from the Ehrenfest theorems
It was established above that the Ehrenfest theorems are consequences of the Schrödinger equation. However, the converse is also true: the Schrödinger equation can be inferred from the Ehrenfest theorems. We begin from
Application of the product rule leads to
Here, apply Stone's theorem, using to denote the quantum generator of time translation. The next step is to show that this is the same as the Hamiltonian operator used in quantum mechanics. Stone's theorem implies
where was introduced as a normalization constant to the balance dimensionality. Since these identities must be valid for any initial state, the averaging can be dropped and the system of commutator equations for are derived:
Assuming that observables of the coordinate and momentum obey the canonical commutation relation . Setting , the commutator equations can be converted into the differential equations
whose solution is the familiar quantum Hamiltonian
Whence, the Schrödinger equation was derived from the Ehrenfest theorems by assuming the canonical commutation relation between the coordinate and momentum. If one assumes that the coordinate and momentum commute, the same computational method leads to the Koopman–von Neumann classical mechanics, which is the Hilbert space formulation of classical mechanics. Therefore, this derivation as well as the derivation of the Koopman–von Neumann mechanics, shows that the essential difference between quantum and classical mechanics reduces to the value of the commutator .
The implications of the Ehrenfest theorem for systems with classically chaotic dynamics are discussed at Scholarpedia article Ehrenfest time and chaos. Due to exponential instability of classical trajectories the Ehrenfest time, on which there is a complete correspondence between quantum and classical evolution, is shown to be logarithmically short being proportional to a logarithm of typical quantum number. For the case of integrable dynamics this time scale is much larger being proportional to a certain power of quantum number.
| Physical sciences | Quantum mechanics | Physics |
2588077 | https://en.wikipedia.org/wiki/Handrail | Handrail | A handrail is a rail that is designed to be grasped by the hand so as to provide safety or support. In Britain, handrails are referred to as banisters. Handrails are commonly used while ascending or descending stairways and escalators in order to prevent injurious falls, and to provide bodily support in bathrooms or similar areas. Handrails are typically supported by balusters or attached to walls.
Similar items not covered in this article include bathroom handrails—which help to prevent falls on slippery, wet floors—other grab bars, used, for instance, in ships' galleys, and barres, which serve as training aids for ballet dancers. Guard rails and balustrades line drop-offs and other dangerous areas, keeping people and vehicles out.
British specifications
British Standard and British Standard Code of Practice are harmonized to European Normal (EN) series. Handrail height is set between .
US specifications
Dimensions
Various model codes—The International Code Council (ICC) and National Fire Protection Association (NFPA)—and accessibility standards—ANSI A117.1 and the Americans With Disabilities Act Standards for Accessible Design (ADASAD)—refer to handrail dimensions. Current versions of these codes and standards now agree that handrail is defined as either a circular cross section with an outside diameter of minimum and maximum or a non-circular cross section with a perimeter dimension of minimum and maximum and a cross section dimension of ) maximum. In addition, the International Residential Code (IRC) includes a definition of a "Type II" handrail that allows for handrail with a perimeter dimension greater than .
The IRC and residential portion of the 2009 IBC define Type II handrail as follows:
Handrails are located at a height between . In areas where children are the principal users of a building or facility, the 2010 ADASAD recommends that a second set of handrails at a maximum height of measured to the top of the gripping surface from the ramp surface or stair nosing can assist in preventing accidents.
Clearance
The distance between the wall and handrail gripping surface is also governed by local code with the most common requirement being minimum. The National Fire Protection Agency (NFPA) and the Occupational Safety and Health Administration (OSHA) require that the distance between the wall and handrail be a minimum of .
The 1992 Americans With Disabilities Act Accessibility Guidelines (ADAAG) stated that there was to be an absolute dimension of between a handrail and a wall. This was actually a "grab bar" dimension which was part of the 1986 ANSI A117.1. ANSI changed the notation to minimum in 1990. This was not corrected in 2010 with the approval of the new ADASAD which now calls for a minimum clearance.
Codes also generally require that there be a clearance between the underside of the handrail and any obstruction—including the horizontal bracket arm. There is an allowance however for variations in the handrail size—for every of additional perimeter dimension over , may be subtracted from the clearance requirement.
Strength
Handrails are to support a continuous load of or a concentrated load of applied at the top of the handrail. As handrails come in different materials, the strengths can vary. From timber to stainless steel, it is best to pick a handrail that will be right for the area. Stainless steel will be stronger and more durable outside, whereas timber can be just as strong but less durable outside.
ADA height notations
ADA Handrail height requirements are issued to create a safe enjoyable space to enjoy for everyone. The ADA height of handrail requirements that will primarily be used by children have their own unique requirements.
The top of gripping surfaces of handrails shall be 34 inches minimum and 38 inches maximum vertically above walking surfaces, stair nosing, and ramp surfaces. Handrails shall be at a consistent height above walking surfaces, stair nosing, and ramp surfaces.
When children are the principal users in a building or facility (e.g., elementary schools), a second set of handrails at an appropriate height can assist them and aid in preventing accidents. A maximum height of 28 inches measured to the top of the gripping surface from the ramp surface or stair nosing. Sufficient vertical clearance between upper and lower handrails, 9 inches minimum, should be provided to help prevent entrapment.
Adult requirements
Top of gripping surfaces of handrails shall be between vertically above walking surfaces, stair nosings, and ramp surfaces. Handrails shall be at a consistent height above walking surfaces, stair nosings, and ramp surfaces.
Child recommendation
When children are the principal users in a building or facility (e.g., elementary schools), a second set of handrails at an appropriate height can assist them and aid in preventing accidents. A maximum height of measured to the top of the gripping surface from the ramp surface or stair nosing is recommended for handrails designed for children. Sufficient vertical clearance between upper and lower handrails, minimum, should be provided to help prevent entrapment.
Types of handrails
Handrails are available in several different varieties including wooden/timber, stainless steel, brass, or aluminium. Some varieties are more hard-wearing than others, and the cost can vary depending on the style you choose. You can also choose various fixings to allow for a more modern/sleek look.
| Technology | Architectural elements | null |
2588079 | https://en.wikipedia.org/wiki/Guard%20rail | Guard rail | Guard rails, guardrails, railings or protective guarding, in general, are a boundary feature and may be a means to prevent or deter access to dangerous or off-limits areas while allowing light and visibility in a greater way than a fence. Common shapes are flat, rounded edge, and tubular in horizontal railings, whereas tetraform spear-headed or ball-finialled are most common in vertical railings around homes. Park and garden railings commonly in metalworking feature swirls, leaves, plate metal areas and/or motifs particularly on and beside gates.
High security railings (particularly if in flat metal then a type of palisade) may instead feature jagged points and most metals are well-suited to anti-climb paint.
A handrail is less restrictive on its own than a guard rail and provides support.
Guardrails also apply in a technology context.
Public safety
Many public spaces are fitted with guard rails as a means of protection against accidental falls. Any abrupt change in elevation where the higher portion is accessible makes a fall possible. Due to this responsibility and liability, rails are placed to protect people using the premises. Guardrails in the US are generally required by code where there is a drop of or more.
Examples of this are both architectural and environmental. Environmental guard rails are placed along hiking trails where adjacent terrain is steep. Railings may also be located at scenic overlooks.
Guard rails in buildings can be numerous, and are required by building codes in many circumstances. Handrails along stairways may be supported by balusters forming a balustrade, and catwalks (a type of footbridge) and balconies are also lined with them. An example of a common residential guard rail (US) handrail (Brit.) is a wood railing around a deck or patio. In the US this is typically built on-site from pressure treated lumber thus featuring a simplistic design of vertical baluster spaced every demonstrating compliance with Building Codes (Standards).
Cable railings typically use stainless steel cables strung horizontally. Glass balusters and glass panels open the view while still providing safety, as at the Grand Canyon Skywalk. With the increasing popularity of composite lumber for decking, manufacturers, like TimberTech are providing composite railing components. Wrought iron is another choice that is traditional and sturdy. Decorative examples are considered ironwork.
Building codes also require that no opening in a guard be of a size such that a sphere may pass. There are three exceptions according to the 2003 International Building Code Section 1012.3 which allow openings to not exceed depending on occupancy groups or special areas.
A major architect imaginatively used handrails representing social stability, Alvar Aalto. The guard rails of an observation tower such as the Space Needle or Eiffel Tower become exaggerated to the point of becoming a fence or cage. This is also done on bridges and overpasses to prevent accidents and suicides.
Facility safety
The majority of safety guardrails used in industrial workplaces are made from fabricated steel. Steel guardrail was originally developed by Armco (The American Rolling Mill Company) in 1933 as highway guardrail but is often used in the factories and warehouses of the industrial sector, despite not being intended for this application. Simultaneously, another version of steel guardrail was developed by Kee Clamp Steel for restraining cattle in the agriculture industry, this was released in 1934 and, like Armco guardrail, is still often used in industrial environments. Box beam steel and I-beam steel are other examples of fabricated steel commonly used to make safety guardrails.
As governments around the world courted the voting power of working people and trade unions in the 1980s, the health and safety rights of workers became of greater importance. This set the governmental procedures in motion that would see a flurry of regulations and legislation around workplace safety being introduced in industrialized countries in the 1990s. In the US and UK, these introductions, combined with the demonstrable success of the Occupational Safety and Health Act 1970 (US), and Health and Safety at Work Act 1974 (UK), led to workplace safety being taken more seriously at industrial facilities. Businesses worldwide began to see the value of effective workplace safety, in both the direct commercial sense of protecting assets, but also in the improved productivity levels of a protected workforce. Part of this increase in desire for industrial workplace protection, was a growth in demand for safety guardrails.
In the US, under OSHA Standard 1910.28(b)(15), employees who work on surfaces that are or higher off of the ground must have personal fall protection systems in place, such as handrails or guardrails. Although the OSHA standard calls for a guardrail or guardrail system to protect workers on elevated work areas, current industry terminology would refer to that type of safety system as a handrail system or safety rail system.
The term "guardrail" as used in industry and distribution facilities refers to floor mounted systems consisting of horizontal rails attached to upright posts that resemble outdoor highway guardrail. The facility safety guardrails control in-plant vehicular traffic and protect areas where errant vehicle contact may cause harm to property or personnel. Common uses for guardrail systems may be along an aisle or building wall. Often guardrails are placed around equipment installations to prevent inadvertent contact from lift trucks. The guardrail provides a safety barrier preventing lift trucks or other in-plant vehicles from traveling into equipment, building walls, or personnel.
There are different types of "industrial" or "facility safety" guardrail systems, each with its own advantages.
Types of guardrail
Steel Ribbed Rail
The most common industrial or facility safety guardrail systems are constructed of steel, where the upright posts are made of heavy wall steel tubing-either round or square, with heavy gauge ribbed steel rails mechanically attached to the uprights by bolts or other fasteners. The uprights may be welded to baseplates and anchored to the floor, or set and concreted into the floor. In industrial and distribution facilities the steel guardrail systems provide solid protection for property and personnel by restricting and controlling where in-plant vehicular traffic may operate. These guardrail systems also serve as a warning for people to be aware of hazardous in facility vehicular traffic. Some facility safety steel guardrail systems utilize one horizontal rail, where others employ two or three rails. While a single rail safety guardrail located some above floor level may be sufficient to control vehicle traffic, it could present a trip hazard to pedestrians since it is not as noticeable as a double or triple guardrail system where the top rail is some above the floor level.
Manufacturers produce several grades of steel safety guardrail. Each grade is suited to a different application which may be determined by size and type of vehicles used in the facility, the volume of traffic to be controlled, or the value/risk associated with the areas being guarded.
Polymer Guardrail
In 1992 a yellow polymer cover sleeve for steel guardrails was introduced by a Detroit-based company which enhanced the visibility of steel safety guardrails in industrial environments and removed the need for continual repainting. This polymer cover was the beginning of polymer guardrail in the industrial environment. At the turn of the 21st century, several companies developed safety products constructed with polymer designed specifically for industrial environments and in 2001 a dedicated industrial alternative to steel guardrail was invented by a British-based company, and thus the first fixed flexible polymer safety guardrail was introduced to the market. In the early 2000s, an Italian company added a version of a polymer bumper providing small bumpers for walls and equipment in the food industry. A Belgian company also introduced a flexible barrier in 2010 and in 2014 a US based company introduced a hybrid polymer-steel guardrail for industrial environments.
There are many types of polymer used in manufacturing safety guardrail. Polymer grains may be blended as part of the production process in order to enhance the natural properties of polymer in different ways. The most common types of polymer used are: polyethylene, polypropylene, and polyvinyl chloride (PVC). Using these types of polymer, there are three basic design philosophies for polymer guardrails:
Impact Dispersal, which applies multi-component engineering to enable guardrails to bend on impact, absorbing and dispersing energy, before returning to their original shape.
Hybrid, which features a steel core encased in polymer.
Central Fix, which utilizes a solid polymer post fixed into the ground with cement and steel subterranean fixings.
The Polymer Impact Dispersal Safety Guardrail design offers an advantage over steel guardrails. Steel guardrails, if impacted with sufficient energy, will permanently deform and require repair or replacement to the impacting vehicle, the guardrail itself, and even the floor substrate. Where the Impact Dispersal Safety Guardrail design allows the guardrail to bend on impact, disperse the energy, and return to its original shape. Resulting in no damage to the guardrail or the impacting vehicle.
Test Standards for Industrial or Facility Safety Guardrail
Until 2017 there were no accepted standards for how safety guardrail tests were conducted, or how the results were communicated. In 2017, the BSI (British Standards Institution) published the publicly available specification, PAS 13:2017 Code of Practice for safety barriers used in traffic management within workplace environments with test methods for safety barrier impact resilience (commonly referred to as PAS13). This outlines test method guidelines for comparing like-for-like guardrail products, as well as being the current best practice traffic management procedures for a workplace and providing a standard for the safety guardrails within them.
In the US, there is no ANSI (American National Standards Institute) Standard for testing guardrails. Steel guardrail system manufacturers test their systems to withstand impacts from a load moving at , whereas manufacturers of polymer safety guardrail systems rate and test their products at varying impact levels due to variation in guardrail product systems they offer, which can range from simple low impact energy pedestrian guardrail systems to high impact energy industrial traffic guardrail systems. In order to standardize testing for both steel and polymer guardrails, the Protective Guarding Manufacturers Association (ProGMA) industry group of the Material Handling Industry (MHI) trade association is currently working with ANSI to develop an industry endorsed testing and reporting standard.
Automotive safety
In traffic engineering, a highway guardrail may prevent an errant vehicle from hitting roadside obstacles, which may be either man-made (sign structures, culvert inlets, utility poles) or natural (trees, rock croppings), running off the road and going down a steep embankment, or veering off the roadway into oncoming traffic (commonly referred to as a median barrier). Roadside obstacles are typically referred to as fixed objects. A secondary objective is keeping the vehicle upright while deflected along the guardrail. Variables such as motorist speed and vehicle orientation when striking the guardrail are crucial factors in the effectiveness of guardrail performance.
The most common type of guardrail in use today is the Blocked-Out W-beam (Strong Post). The strong-post W-beam guardrail consists of wood posts and wood blockouts or steel posts with wood or plastic blockouts. The wood or plastic blockouts reduce or minimize a vehicle snagging on the posts upon impact. In addition, a blockout may be used to increase the offset of guardrail with an obstacle such as a curb. The posts' primary purpose is to maintain the height of the guardrail during the initial stages of post deflection. Maintaining guardrail height also reduces the potential for a vehicle to vault over the guardrail upon initial impact. The posts also play a role in the amount of resistance and deflection a guardrail may experience during impact. Resistance in a strong post system results from a combination of tensile and flexural stiffness of the rail and the bending and shearing resistance of the posts.
One of the main concerns with strong-post W-beam guardrail has been the ability of the system to contain and redirect modern vehicles that have a higher center of gravity along with the increased weight of those vehicles. For instance the optimum height of the strong-post W-beam for a car might not keep a truck from toppling over it, while a motorbike might slip under a higher rail. To address these concerns, significant research and development of a system that could contain and redirect vehicles of varying weights and heights was developed and crash tested (both controlled and simulated). As a result, the Midwest Guardrail System (MGS) was developed and successfully crash tested per NCHRP Report 350 TL-3 criteria. The MGS has a higher mounting height, uses larger posts and blockouts as compared to the strong-post W-beam guardrail. One other significant difference is that MGS rail splices occur at mid-span compared to at the posts like the strong-post W-beam guardrail.
In most cases, a guardrail would not be able to withstand the impact of a vehicle just by the strength of the individual posts in the area hit by the vehicle. Guardrailing functions as a system with the guardrail, posts, connection of the rail to the posts and to each other, and the end anchors (or terminals) all playing an integral role in how the guardrail will function upon impact. Soil conditions, height of rail, presence of curb or dike, weight of impacting vehicle, distance from back of post to hinge point and depth of post within soil can all determine how well the system will function upon impact.
A guardrail is effectively one strong band that transfers the force of the vehicle to the rail elements, posts, and end terminals or anchors. A run of guardrail must be anchored at each terminating end either by transitioning the rail into a fixed anchor such as a bridge rail or with an end terminal or end anchor placed in the ground or within an embankment. Newer concrete barriers, while usually strong enough to withstand direct hits by cars, still work on a similar principle in deflecting heavier vehicles such as trucks.
Guard rails are intended to deflect. The amount of deflection is dependent on a number of factors, some of which include the type and weight of impacting vehicle, height the guardrail is placed, type of soil the posts may be embedded within, length of embedment of the posts, and distance of the hinge point to the face of the guardrail are just a few. A guardrail that deflects significantly can causes pocketing, which has the potential to snag a vehicle, which may cause it to flip, roll, or cause the rail to fail entirely, allowing a vehicle to penetrate the guardrail.
Modern installations of guardrails are designed to allow the guiderail to deform under the load of the crash, and safely redirecting a vehicle back onto the roadway at a somewhat shallow angle. It is important that the approach grades to a guardrail system be very flat (typically 10:1 or flatter) and that grades and fixed objects behind guardrail be placed at a distance so that it will not affect the performance of the guardrail upon impact and deflection.
Absorption is when the force of impact is directly transferred between the vehicle and guardrail, which may cause the end to puncture the vehicle. This is most common where a "whale tail" or blunt end treatment exists. To mitigate this a number of guiderail end treatments exist such as "extruder end treatments", "eccentric loaders" and "driveway wrap treatments" which result in blunt ends rarely being left exposed in modern installations.
Lastly, a vehicle can become airborne upon striking a guardrail with a buried end treatment if the slope to which the end anchor is buried is relatively flat (3:1 or flatter), which may negate the purpose of the guardrail, if the vehicle continues beyond the guardrail and strikes the object the guardrail was protecting. Additionally, an airborne vehicle is likely to collide in a manner that the vehicle was not designed for, increasing the risk of failure in the vehicle's collision safety systems. A guardrail will have some give and deflection upon impact. The amount of deflection depends on many factors of which speed and weight of vehicle, type of guardrail installed, height of rail, length of posts, soil conditions and a number of other factors can all play a role. Guardrailing must be installed so that it is not so rigid that the rail will fail upon impact or the posts will snap off at the point where they are embedded within the ground.
Transportation engineers limit the amount of guardrail placed as much as possible, as guardrails should only be placed when the roadside conditions pose a greater threat than the guardrail itself. In fact, in the hierarchy of five roadside safety treatments, shielding with guardrails ranks fourth. Therefore, while guardrails are often added as a retrofit to existing roads, newer roads are designed to minimize roadside obstacles, whether that may include aligning a road on a smoother curve or filling in a ravine which would eliminate the need for guardrail altogether. In addition to new research into end treatments, public awareness among both drivers and engineers has been gradually reducing injuries and fatalities due to guardrails. Though they have usually prevented far more serious accidents, guardrails are considered roadside obstacles as well and transportation engineers must weigh whether placing a guardrail will reduce the severity of an impact as compared to what may be impacted if the guardrail were not placed. In general, the minimum length of guardrail with an end anchor at the trailing end and an end terminal on the approach end will be in length. An example would be where an overhead roadside sign structure is placed within what is considered the clear recovery zone—an engineer would need to determine that the structure has a potential to be impacted and the impact of a vehicle with that structure would be much more severe than impacting the guardrail.
There are four general types of guardrail, ranging from weakest and inexpensive to strongest and expensive; cable and wood posts, steel and wood/metal posts, steel box-beam, and concrete barriers. While cheaper guardrail is the weakest, often being destroyed from the impact of a light vehicle, it is inexpensive and quick to repair, so this is frequently used in low-traffic rural areas. On the other hand, concrete barriers can usually withstand direct hits from a larger variety of vehicle types, making them well suited for use on high volume routes such as freeways or ramps with sharp curves. While rarely damaged, they would be considerably more expensive and time-consuming to repair. Concrete barriers are frequently installed in the median, being expected to withstand frequent impacts from both sides, while the shoulders of the road often have cheaper guardrail. Although the use of concrete barriers on the right side of highway is becoming a much more frequent occurrence in areas where guardrail may be sustaining frequent impacts and the ability for maintenance repairs may be restricted by the general area or work windows due to high traffic volumes for most of the day.
Traffic dangers
In cities occasionally pedestrian railings (and barriers) are installed at the immediate side of a roads. However, cyclists have died when crushed against them by motor vehicles. Close "safety barriers" to roads have been found to increase the chances of injury to pedestrians for a number of reasons including increasing inattention of drivers and pedestrians. For these reasons some councils in the United Kingdom have removed their pedestrian railings. This was after London's Royal Borough of Kensington and Chelsea did so and found that the rate of injury to pedestrians decreased three times faster than elsewhere in the city.
The removal of barriers divorcing wheeled traffic from pedestrians is an essential of the shared space paradigm and philosophy in urban design. Security barriers have been introduced on a number of the world's major bridges and around key areas of congregation.
Railways
Railway trackage has guard rails (aka check rails) to guide wheels through possible catch points on turnouts or diamonds. Similarly, guard rails may be installed inside the innermost running rails on very sharp curves. The other most common usage is to prevent damage to other structures, especially bridges, in a derailment.
Technology
In a technology context, a guardrail is an artifact that defines the boundaries in which technology change can be executed in a manner that is aligned with organisational strategy, risk, architecture, operational and cyber security requirements.
Examples of technology guardrails are:
Principles
Policies
Strategies
Technical Standards
Patterns
Guidelines
Reference Architectures (conceptual, logical, physical)
Each of these guardrails constrain what technology teams can do within approved boundaries (the scope of the guardrail). For example, a technology standard may define a certain database as the standard to use within an organisation. The approved standard would have a commercial agreement, operational and delivery capability established, functional and non-functional fit for purpose assessments undertaken, and a cyber security assessment completed. Technology delivery teams are constrained by the Guardrail to use that particular database technology, but can do so confidently for the scenarios in which usage has been defined and approved. If a technology delivery team chose a different database technology that is not defined within the technology standards, they could be introducing commercial, organisational or operational risk as the pre-requisite due diligence has not been undertaken.
| Technology | Architectural elements | null |
2589115 | https://en.wikipedia.org/wiki/NAND%20gate | NAND gate | In digital electronics, a NAND gate (NOT-AND) is a logic gate which produces an output which is false only if all its inputs are true; thus its output is complement to that of an AND gate. A LOW (0) output results only if all the inputs to the gate are HIGH (1); if any input is LOW (0), a HIGH (1) output results. A NAND gate is made using transistors and junction diodes. By De Morgan's laws, a two-input NAND gate's logic may be expressed as , making a NAND gate equivalent to inverters followed by an OR gate.
The NAND gate is significant because any Boolean function can be implemented by using a combination of NAND gates. This property is called "functional completeness". It shares this property with the NOR gate. Digital systems employing certain logic circuits take advantage of NAND's functional completeness.
NAND gates with two or more inputs are available as integrated circuits in transistor–transistor logic, CMOS, and other logic families.
Symbols
There are three symbols for NAND gates: the MIL/ANSI symbol, the IEC symbol and the deprecated DIN symbol sometimes found on old schematics. The ANSI symbol for the NAND gate is a standard AND gate with an inversion bubble connected.
Logic
The function is logically equivalent to
One way of expressing A NAND B is , where the symbol signifies AND and the bar signifies the negation of the expression under it: in essence, simply .
Implementations
The basic implementations can be understood from the image on the left below: If either of the switches S1 or S2 is open, the pull-up resistor R will set the output signal Q to 1 (high). If S1 and S2 are both closed, the pull-up resistor will be overridden by the switches, and the output will be 0 (low).
In the depletion-load NMOS logic realization in the middle below, the switches are the transistors T2 and T3, and the transistor T1 fulfills the function of the pull-up resistor.
In the CMOS realization on the right below, the switches are the n-type transistors T3 and T4, and the pull-up resistor is made up of the p-type transistors T1 and T2, which form the complement of transistors T3 and T4.
In CMOS, NAND gates are more efficient than NOR gates. This is due to the faster charge mobility in n-MOSFETs compared to p-MOSFETs, so that the parallel connection of two p-MOSFETs (T1 and T2) realised in the NAND gate is more favourable than their series connection in the NOR gate. For this reason, NAND gates are generally preferred over NOR gates in CMOS circuits.
Hardware design and pinout
NAND gates are basic logic gates, and as such they are recognised in TTL and CMOS ICs.
The standard, 4000 series, CMOS IC is the 4011, which includes four independent, two-input, NAND gates. These devices are available from many semiconductor manufacturers. These are usually available in both through-hole DIL and SOIC formats. Datasheets are readily available in most datasheet databases.
The standard two-, three-, four- and eight-input NAND gates are available:
CMOS
4011: Quad two-input NAND gate
4023: Triple three-input NAND gate
4012: Dual four-input NAND gate
4068: Mono eight-input NAND gate
TTL
7400: Quad two-input NAND gate
7410: Triple three-input NAND gate
7420: Dual four-input NAND gate
7430: Mono eight-input NAND gate
Functional completeness
The NAND gate has the property of functional completeness, which it shares with the NOR gate. That is, any other logic function (AND, OR, etc.) can be implemented using only NAND gates. An entire processor can be created using NAND gates alone. In TTL ICs using multiple-emitter transistors, it also requires fewer transistors than a NOR gate.
As NOR gates are also functionally complete, if no specific NAND gates are available, one can be made from NOR gates using NOR logic.
| Technology | Digital logic | null |
6214456 | https://en.wikipedia.org/wiki/Evolution%20of%20spiders | Evolution of spiders | Spiders have existed since at least 380 million years. The group's origins lie within an arachnid sub-group defined by the presence of book lungs (the tetrapulmonates); the arachnids as a whole evolved from aquatic chelicerate ancestors. More than 45,000 extant species have been described, organised taxonomically in 3,958 genera and 114 families. There may be more than 120,000 species. Fossil diversity rates make up a larger proportion than extant diversity would suggest with 1,593 arachnid species described out of 1,952 recognized chelicerates. Both extant and fossil species are described annually by researchers in the field. Major developments in spider evolution include the development of spinnerets and silk secretion.
Early spider-like arachnids
Among the oldest known land arthropods are Trigonotarbids, members of an extinct order of spider-like arachnids.
Trigonotarbids share many superficial characteristics with spiders, including a terrestrial lifestyle, respiration through book lungs, and walking on eight legs, with a pair of leg-like pedipalps near the mouth and mouth parts. They lacked the ability to spin silk: there is no evidence for either spigots or spinnerets within the group. An unpublished fossil exists which has distinct microtubercles on its hind legs, akin to those used by spiders to direct and manipulate their silk, but given the lack of any structures associated silk production, it seems unlikely the structures were associated with silk.
Trigonotarbids are not true spiders, and the trigonotarbids have no living descendants.
Emergence of true spiders
According to a 2020 study using a molecular clock calibrated with 27 chelicerate fossils, spiders most likely diverged from other chelicerates between 375 and 328 million years ago.
At one stage, Attercopus was claimed as the oldest fossil spider which lived during the Devonian. Attercopus was placed as the sister-taxon to all living spiders, but has now been reinterpreted as a member of a separate, extinct order Uraraneida which could produce silk, but did not have true spinnerets. The discovery of Chimerarachne in early Late Cretaceous (Cenomanian) aged Burmese amber has also demonstrated that taxa existed until the Cretaceous that had both spinnerets, and a whip-like telson.
The oldest reported spiders date to the Carboniferous Period, or about . Most of these early segmented fossil spiders from the Coal Measures of Europe and North America probably belonged to the Mesothelae, or something very similar, a group of spiders with the spinnerets placed underneath the middle of the abdomen, rather than at the end as in modern spiders. They were probably ground-dwelling predators, living in the giant clubmoss and fern forests of the mid-late Palaeozoic, where they were presumably predators of other primitive arthropods. Silk may have been used simply as a protective covering for the eggs, a lining for a retreat hole, and later perhaps for simple ground sheet web and trapdoor construction. They co-existed with a range of spider-like forms which had some, but not all, the characters associated with the true spiders.
As plant and insect life diversified so also did the spider's use of silk. Spiders with spinnerets at the end of the abdomen (Mygalomorphae and Araneomorphae) appeared more than , presumably promoting the development of more elaborate sheet and maze webs for prey capture both on ground and foliage, as well as the development of the safety dragline. The oldest mygalomorph, Rosamygale, was described from the Triassic of France. Megarachne servinei from the Permo-Carboniferous was once thought to be a giant mygalomorph spider and, with its body length of and leg span of above , the largest known spider ever to have lived on Earth, but subsequent examination by an expert revealed that it was actually a relatively small sea scorpion.
By the Jurassic period, the sophisticated aerial webs of the orb-weaver spiders had already developed to take advantage of the rapidly diversifying groups of insects. A spider web preserved in amber, thought to be 110 million years old, shows evidence of a perfect "orb" web, the most famous, circular kind one thinks of when imagining spider webs. An examination of the drift of those genes thought to be used to produce the web-spinning behavior suggests that orb spinning was in an advanced state as many as . One of these, the araneid Mongolarachne jurassica, from about , recorded from Daohuogo, Inner Mongolia in China, is the largest known fossil of a spider.
The 110 million year-old amber-preserved web is also the oldest to show trapped insects, containing a beetle, a mite, a wasp's leg, and a fly. The ability to weave orb webs is thought to have been "lost", and sometimes even re-evolved or evolved separately, in different species of spiders since its first appearance.
Around half of modern spider species belong to the RTA clade, a group of spiders linked by the shared morphological trait of the retrolateral tibial apophysis (RTA) on the male pedipalp. Despite their modern diversity, there is no unambiguous evidence of the clade from the Mesozoic, though molecular clocks suggest that diversification of the group began in the Late Cretaceous. There appears to be a faunal turnover in the Cretaceous-Cenozoic interval, with the Cretaceous dominated by Synspermiata and Palpimanoidea, as well as enigmatic extinct families like the lagonomegopids, while the Cenozoic is dominated by RTA clade and araneoid spiders.
| Biology and health sciences | Basics_4 | Biology |
6221709 | https://en.wikipedia.org/wiki/Cold%20wave | Cold wave | A cold wave (known in some regions as a cold snap, cold spell or Arctic Snap) is a weather phenomenon that is distinguished by a cooling of the air. Specifically, as used by the U.S. National Weather Service, a cold wave is a rapid fall in temperature within a 24-hour period requiring substantially increased protection to agriculture, industry, commerce, and social activities. The precise criteria for a cold wave are the rate at which the temperature falls, and the minimum to which it falls. This minimum temperature is dependent on the geographical region and time of year.
In the United States, a cold spell is defined as the national average high temperature dropping below . A cold wave of sufficient magnitude and duration may be classified as a cold air outbreak (CAO).
Effects
A cold wave can cause death and injury to livestock and wildlife. Exposure to cold mandates greater caloric intake for all animals, including humans, and if a cold wave is accompanied by heavy and persistent snow, grazing animals may be unable to reach needed food and die of hypothermia or starvation. They often necessitate the purchase of foodstuffs to feed livestock at considerable cost to farmers.
Cold spells are associated with increased mortality rates in populations around the world. Both cold waves and heat waves cause deaths, though different groups of people may be susceptible to different weather events. More temperature-attributable deaths occur during a cold wave than in a heat wave, though the mortality rate is higher in undeveloped regions of the world. Extreme winter cold often causes poorly insulated water pipelines and mains to freeze. Even some poorly protected indoor plumbing ruptures as water expands within them, causing much damage to property and costly insurance claims. Demand for electrical power and fuels rises dramatically during such times, even though the generation of electrical power may fail due to the freezing of water necessary for the generation of hydroelectricity. Some metals may become brittle at low temperatures. Motor vehicles may fail when antifreeze fails or motor oil gels, producing a failure of the transportation system.
Fires become even more of a hazard during extreme cold. Water mains may break and water supplies may become unreliable, making firefighting more difficult. The air during a cold wave is typically denser and thus contains more oxygen, so when air that a fire draws in becomes unusually cold it is likely to cause a more intense fire. However, snow may stop spreading of fires, especially wildfires.
Winter cold waves that are not considered cold in some areas, but cause temperatures significantly below average for an area, are also destructive. Areas with subtropical climates may recognize a cold wave at higher temperatures than other, colder areas of the globe. The cold wave may be recognized at barely freezing temperatures, as these are still unusually cold for the region, and plant and animal life will be less tolerant of such cold. The same winter temperatures that one associates with the norm for Colorado, Ohio, or Bavaria are catastrophic to crops in places like Florida, California, or parts of South America that grow fruit and vegetables in winter.
Cold waves that bring unexpected freezes and frosts during the growing season in mid-latitude zones can kill plants during the early and most vulnerable stages of growth, resulting in crop failure as plants are killed before they can be harvested economically. Such cold waves have caused famines. At times as deadly to plants as drought, cold waves can leave land in danger of later brush and forest fires that consume dead biomass. One extreme was the so-called Year Without a Summer of 1816, one of several years during the 1810s in which numerous crops failed during freakish summer cold snaps after volcanic eruptions that reduced incoming sunlight.
Recent research suggests a possible link between cold waves in North America and extratropical cyclogenesis over the East Atlantic. These may be connected by large-scale atmospheric circulation patterns. Examples include Rossby wave propagation from the North Pacific or an upper-level anticyclone west of Greenland. In Europe, the advection of cold air masses from the northeast emerges as a potential precursor signal for the majority of cold waves, significantly affecting the energy, health, and agricultural sectors of the continent.
Countermeasures
In some places, such as Siberia, extreme cold requires that fuel-powered machinery intended to be used occasionally must be run continually. Internal plumbing can be wrapped, and persons can often run water continuously through pipes. Energy conservation, difficult as it is in a cold wave, may require such measures as collecting people (especially the poor and elderly) in communal shelters. Even the homeless may be arrested and taken to shelters, only to be released when the hazard abates. Hospitals can prepare for the admission of victims of frostbite and hypothermia; schools and other public buildings can be converted into shelters.
People can stock up on food, water, and other necessities before a cold wave. Some may even choose to migrate to places of milder climates, at least during the winter. Suitable stocks of forage can be secured before cold waves for livestock, and livestock in vulnerable areas might be shipped from affected areas or even slaughtered. Smudge pots can bring smoke that prevents hard freezes on a farm or grove. Vulnerable crops may be sprayed with water that will paradoxically protect the plants by freezing and absorbing the cold from surrounding air.
Most people can dress appropriately and can layer their clothing should they need to go outside or should their heating fail. They can also stock candles, matches, flashlights, and portable fuel for cooking and wood for fireplaces or wood stoves, as necessary. However, caution should be taken as the use of charcoal fires for cooking or heating within an enclosed dwelling is extremely dangerous due to carbon monoxide poisoning. Adults must remain aware of the exposure that children and the elderly have to cold.
Historical cold waves
17th century cold waves (1601–1700)
Great Frost of 1683–84, the worst frost in England in its history. The cold caused the entire River Thames to freeze up to a depth of . The frost enabled the River Thames Frost Fair. See Maunder Minimum
Europe winter of 1694–1695. See Maunder Minimum
18th century cold waves (1701–1800)
Great Frost of 1709, the coldest winter in Europe ever recorded.
19th century cold waves (1801–1900)
1835
Eastern cold wave of January and February 1835. First of three historic U.S. cold waves to hit during the 19th century (1835, 1857, 1899). In January, mercury thermometers froze throughout the Northeast. Mercury froze at in Bangor and Bath, Maine and Montpelier and White River, Vermont. In Connecticut, Hartford hit and New Haven , and in Massachusetts, Williamstown hit and Pittsfield , all low temperature marks that have never been matched since. In February, Savannah, Georgia the temperature hit , colder than would be reached during the 1899 cold wave later in the century and Charleston South Carolina hit .
1836
Last reported snowfall in Sydney, Australia occurred on 28 June of that year. British settlers in Hyde Park woke up to snow "nearly deep", with the meteorological table in The Sydney Herald recording that on the morning of the snow the temperature dropped to .
1857
New England Cold wave of 1857. January 1857 was the coldest month ever recorded in New England. Average month temperatures of in New Haven, in Boston, and in New York City remain coldest months on record in those cities. The worst of the cold descended on New England on January 22 with January 23 being one of the coldest days known in the region. In Bath, Maine a temperature reading of and in Franconia, New Hampshire were recorded. In Norwich, Vermont was recorded. Boston suburbs of Malden and West Newton recorded overnight. Boston temperatures for January 23 never rose above all day and Nantucket Island was connected to the mainland by ice. In New York City, Erasmus Hall in Brooklyn reached a high of 0 °F during the day and the Hudson River froze over solidly enough for people to walk across to Hoboken.
1859
January 1859: January 10 – coldest single daytime temperatures ever recorded experienced in New York City and in New England areas. Montreal recorded temperature of at 7 am, some degrees lower than modern Montreal record of in 1933. Toronto recorded on the same day. At the University of Vermont in Burlington, was recorded at 7 am and at 2 pm. In Woodstock, Vermont a temperature of was recorded. Harvard College recorded at 2 pm and the next morning of January 11, the lowest known temperature reading recorded in Boston. Nantucket Island measured , colder than the modern known record. In New York City, recorded temperatures did not go above . In Brooklyn Heights, a recorded reading of at noon and in Eramus Hall in Brooklyn recorded a high of at 7 am and at 9 pm that night. Union Hall in Jamaica Queens recorded at midnight between January 10 and 11. In White Plains, there were readings of at 7 am, at 2 pm, and at 9 pm.
1874–1875
Winter 1874–1875 in Mid-Western United States.
1882–1883
Winter 1882–1883 in United States.
1886–1887
Winter of 1886–87 in the United States Great Plains and Upper Midwest.
1888
1888 US cold wave – A severe cold wave that passed through the Pacific Northwest. It led to a blizzard for the northern Plains and upper Mississippi valley where many children were trapped in schoolhouses where they froze to death.
1893
1893 East Asia Cold Wave – Produced snow in Hong Kong and South China, and freezing temperatures into tropical latitudes.
1893 Eastern United States Cold Wave.
1895
February 1895 United States Cold Wave (Great Freeze) – Damaged citrus crops in Florida. A snowstorm produced unprecedented snowfall amounts along the Gulf Coast, including 22 inches (56 cm) in Houston, TX. Snow fell as far south as Tampico, Mexico, the lowest latitude in North America that snow has been recorded at sea level.
Winter of 1894–95 in the United Kingdom.
1899
February 1899 Cold Wave – Still ranked as number one cold wave outbreak in U.S. history to date.
20th-century cold waves (1901–2000)
1904
The winter of 1904 was the coolest year on record worldwide.
1912
January 1912 cold wave – The severe 1912 United States cold wave caused the longest recorded period of weather below .
1916-1917
Winter of 1916–1917 – the "extended winter" (October to March) of 1916–17 was the coldest on record in the West and Midwest.
1917-1918
Winter of 1917–1918 – The winter was very frigid across the East and created a heating fuel crisis equaled only in January 1977. Severe cold wave in December 1917 and January 1918 in northeast. December 30 set a number of record lows at the time in New York City () and Boston (). Under ideal conditions for radiational cooling, including fresh snow cover and mostly clear skies, the morning of December 30, 1917, was exceptionally cold also in parts of Virginia and West Virginia, with all-time record cold temperatures (that stand until nowadays) recorded in many cities, including at Lewisburg (West Virginia state record), at White Sulphur Springs, West Virginia, at Bluefield, West Virginia, and at Blacksburg and Burke's Garden, Virginia. January 1918 also brought persistent well below average temperatures for many parts of the East and Midwest, with another shot of very cold air in early February (New York City Central Park's high of only on February 5 is a monthly record, while Michigan's Houghton Lake reached a low of on February 1). The Ohio River froze solidly along its entire length.
This cold wave occurred not only in America but also East Asia. In Seoul, the weather was warm in early December 1917, but the temperature didn't go above () during the cold wave from December 15–January 9. In addition, the temperature went lower than () from December 15–31. The cold wave was especially fierce from December 26–27 and the low temperature on December 26 was (), while the high temperature was (). That day Incheon recorded a high temperature of (), this was the 3rd lowest high temperature recorded during winter. Also, the low temperature on December 27 was (), this temperature is 2nd coldest temperature recorded in December in Seoul. (The lowest was () on December 31, 1927). The cold wave continued into January 1918, where it intensified. The low temperature on January 3 was (). The temperature was a slightly warmer on January 14, but the weather was still quite cold, and the low temperatures were at () or lower until January 28. Leading into February, this cold wave dissipated. Seoul was had an average temperature of () in December 1917, the lowest average recorded in December, and recorded an average temperature of () in January 1918, the 6th lowest recorded in January.
The cold wave also impacted North Korea, and the lowest temperature in Pyeongyang was () in December 1917. Also, the average temperature was ().)
1930
A cold wave gripped the western United States in January 1930. Two inches of snow fell in Palm Springs, CA on January 11, one of only two times in the city's history that snow was ever observed.
1932
Major cold outbreaks affected California in January, February and December. Up to two inches of snow fell across the Los Angeles Basin on January 15, and two inches of snow was officially recorded at the Downtown Los Angeles Weather Bureau Office. Snow also fell in San Francisco on three days in December 1932.
1933
1933 Western United States cold wave – The winter of 1932–33 was the second- or third-coldest on record in most of the West (the coldest on record in Arizona) and saw record cold temperatures in Seneca, Oregon (-54 °F/-48 °C), Moran, Wyoming (-66 °F/-54 °C) and Seminole, Texas (-23 °F/-31 °F) between February 7 and 10, when sixty deaths were blamed on extreme cold and ice storms.
1934
February 1934 Cold Wave in New England and Eastern Canada – Longest period of cold weather ever experienced to this point. Average temperatures in upper New England and Eastern Canada were around zero degrees Fahrenheit for most of the month. Lake Ontario was reportedly completely frozen over. Temperatures reached above freezing only on one day in Burlington, VT in February.
1936
1936 North American cold wave – The cold wave of 1936 was the only cold wave of the 1930s to severely impact the United States east of the hundredth meridian. One of the coldest winters in the Great Plains on record. Low temperatures dropped below in Malta, Montana on four separate days and most of Montana averaged 20 degrees below normal for the entire month of February Parshall, North Dakota hit on February 15, still a record. Langdon, North Dakota remained below for 41 straight days from January 11 to February 20, the longest stretch in recorded history for the U.S. outside of Alaska. The cold wave was followed by one of the hottest summers on record, the 1936 North American heat wave.
1937
1937 Western United States cold wave – January 1937 was the coldest month on record in the West and saw snowfall as far south as the hot desert city of Yuma, Arizona, for one of only two occasions on record. California and Nevada saw their lowest temperatures on record: at Boca on January 20 and at San Jacinto on January 8.
1940
January 1940 Southern United States cold wave – Late January saw record-breaking cold and snow across the Southern United States. It was the coldest month there since February 1899.
1941–1942
Winter of 1941–42 in Eastern Europe – The winter of 1941–42 was the coldest of the twentieth century in most of Eastern Europe (e.g. Moscow) and was the last of a succession of abnormally cold winters there that affected the course of World War II.
1947
Winter of 1946–47 in the United Kingdom
February 1947 saw the coldest temperature on record in Canada, at -62.8 °C at Snag, Yukon.
1949
January 1949 Western United States cold wave – The winter of 1948–49 was the coldest since 1891 over the Western United States and saw record snowfall, ice storms as far south as Texas, and constant disruptions to surface transport, along with large losses in livestock and crops. Coldest winter was recorded in many places in California, Nevada, Idaho and Washington state. The cold was also accompanied by severe blizzards which isolated Wyoming ranches and paralyzed the Great Basin region. The U.S. Army ran "Operation Hay Lift" in the region to bring food and hay by plane to isolated ranches in the region. Las Vegas Nevada got a record of snowfall during the month of January. Snow fell in both San Diego and Los Angeles on three days in January 1949. All-time record low of 0 °F in San Antonio, Texas.
1950
1950 Northwest North American cold wave – January 1950 saw unprecedented cold and snowfall in the Pacific Northwest, with normally mild Seattle and Portland, Oregon, both falling below and receiving extremely heavy snow that disrupted transport and schooling as it could not be removed easily. Western Canada saw by far its coldest month on record, leading to severe damage to fruit crops in the Okanagan Valley, the freezing of Okanagan Lake for the only time since 1862, and Calgary's only month where temperatures remained below throughout. Vancouver, British Columbia, had an average temperature of , compared to the average .
1954-1955
Winter of 1954–1955 in East Asia – One of the coldest winters on record across China. Numerous major rivers and lakes froze over across southern China, including the Huai River, Han River, and Dongting Lake.
Between Late July and Early August 1955 the most intense cold wave ever recorded in Brazil was registered. It snowed heavily for over 24 hours in some cities of the south of the country, accumulating over 2 feet (60 cm) in the mountainous regions of Rio Grande do Sul and Santa Catarina states. The cold air reached the Amazon and even crossed the Equator, which is extremely rare. Temperatures plummeted to -10 °C (14 °F) in Bom Jesus, Rio Grande do Sul, and -2 °C (29 °F) in São Paulo (negative temperatures were never again recorded in the city center).
1956
1956 European cold wave – February 1956 was the coldest month of the twentieth century over large areas of Western Europe, with mean temperatures below as far south as Marseilles being utterly unprecedented in records dating back into the eighteenth century.
1962–1963
Winter of 1962–63 in the United Kingdom – The winter of 1962–63 was the coldest for 223 years in England, and the freeze was accompanied by strong easterly winds and the freezing of rivers and streams.
January 1963 cold wave in Mid-Western United States, as well as a brief-but-severe cold spell in the western United States.
1966
1966 Western Canadian cold wave – January 1966 was the coldest January on record in the Yukon and the coldest since 1950 or 1936 in the Prairie Provinces, and the severe cold continued into March, when Winnipeg recorded its most severe winter snowstorm on record.
1968–1969
Winter of 1968–69 in Central Asia – Central Asia and western Siberia saw by far their coldest winter on record in 1968–69, and in Central Asia also their wettest, producing record low temperature, severe blizzards and avalanches, numerous plant deaths and record spring flooding. The cold occasionally swept into East Asia, resulting in record snowstorms and cold in China and Japan.
1969 Northwest North American cold wave – December 1968 and January 1969 saw record cold and snow in the Pacific Northwest and Southern BC. Vancouver, BC recorded its coldest temperature on record of on January 23, 1969, and at the airport on December 29, 1968. Seattle recorded its snowiest winter on record with for the season at Sea-Tac Airport.
1975
July 1975 was a historic month in South America. One of the most intense cold waves of the century climbed through the continent, even crossing the Equator in the Amazon Forest. On July 16, snow fell heavily in Argentina, and on the following day in Paraguay and Southern Brazil. In Curitiba it snowed for around 6 hours, accumulating on the ground, even in the city center. The phenomenon was registered in 5 states, a very rare occurrence. In July 18, the temperature dropped even more. In the state of Paraná, the coffee crops were killed by an episode of black frost (it occurs when the plants' tissues freeze and die), and some cities recorded -10 °C (14 °F), among the lowest temperatures ever recorded in the country. The cold wave reached as far as 10°N before dissipating.
1977
Cold wave of January 1977. Greatest eastern US cold wave of the 20th century. The core of the cold air extended from New Hampshire to Florida and west to Iowa and Missouri. Ohio was at the very center of the cold air mass where every weather station there recorded its coldest month on record. Cincinnati recorded its lowest known temperature of dating back to 1820. The South Carolina state record temperature of was recorded during this cold wave near Long Creek. The wind chill in Minneapolis was on January 28, possibility the lowest ever recorded there up until that point. Snow fell in Miami and Homestead Florida, the farthest south snow was ever recorded in America. President Jimmy Carter walked in his inauguration parade in temperatures below freezing on January 20. Buffalo, New York was hit with its worst blizzard ever during the last week of January where near hurricane-force winds created whiteout conditions for three days. Temperatures in Buffalo were around , wind chills recorded of using the old formula, and the blizzard paralyzed the city with snow drifts of up to .
1978
Cold wave of early 1978 – Produced one of the coldest winters on record in all states east of the Rockies, except Maine.
Europe and Asia, winter of 1978–1979, caused by the Kara Sea 1978 anticyclone. Weather conditions typical for polar regions were detected in Moscow, Leningrad (St. Petersburg), and Sverdlovsk (Yekaterinburg), affecting logistics and the energy industry, and causing fires at the Beloyarskaya Nuclear powerplant. Freeze, but a lack of snow caused winter cereal crop failure throughout 1979.
1979
Cold wave of 1979 – widespread cold across the United States. One of the largest Chicago snowstorms in history at the time, with 21 inches of snowfall in the two-day period, the 1979 Chicago Blizzard occurred during the cold wave in January.
Late 1970s (1977, 1978, 1979) – In the last three years of the 1970s, almost all of the conterminous United States had at least one winter with a memorable cold wave, and the winter of 1978–79 was the coldest on record in the lower 48, with everywhere, except normally frigid upstate Maine, experiencing well below average temperatures.
1981–1982
Winter of 1981/82 in the United Kingdom – This was a significantly colder than average winter. December started off very mild with temperatures up to , but it quickly became very cold and snowy. The night of the 12th–13th is particularly noted for its cold temperatures with many records broken. January 1982 was also a cold and snowy month with records being broken on the 10th in both England and Scotland. England recorded a record low of and down to in Braemar.
January 1982 cold air outbreak – January 1982 was very cold. The 1981 AFC Championship Game, held in Cincinnati, was nicknamed the "Freezer Bowl" due to the temperature at kickoff and wind chill. The Sunday of the following week (January 17, 1982) is also known as Cold Sunday. Chicago's Midway and O'Hare airports record their all-time low temperatures of . Milwaukee, Wisconsin recorded temperatures of on January 17, the lowest in 111 years there. Recorded temperature of in Atlanta and Jackson, Mississippi.
1983
December 1983 Great Plains cold wave – The contiguous US had its coldest ever Christmas in 1983, except for the southwestern US. Severely cold winds blew in from Canada and about 70% of the month was colder than average. Many locations east of the Rockies broke December cold records on Christmas Eve. In addition to cold, the Sioux Falls area had winds bringing wind chills down to . High temperatures did not even reach in northern Illinois during the days before Christmas. Temperatures dropped below on December 15 and remained there for over nine days at Sioux Falls. Minneapolis recorded an average temperature for the month of , the coldest on record. In Chicago, a temperature of and 30 mph winds resulted in a wind chill of ( under the new formula) on Christmas Eve.
1985–1986
1985 Great Western cold air outbreak – February 1985 saw the contiguous U.S.'s second-coldest temperature of in Peter Sinks, Utah. About a month of severe cold affected a large part of the nation. 1985 became the fourth-coldest calendar year on record in the Pacific Northwest.
January 1985 – January 1985 was the coldest January since 1979 in the United Kingdom with significantly below-average temperatures.
January 1985 US cold air outbreak – On January 21, 1985, it was so cold that President Ronald Reagan's inauguration took place in the Capitol Rotunda. In addition to the cold in Washington, D.C., Miami Beach recorded its only frost since records began, lasting for a full three hours. Several other Southern cities set all-time record cold.
Winter of 1985/86 in the United Kingdom – The cold weather started in November 1985 with the month being considerably below average, being the coldest since at least 1925. December 1985 was a milder month and January was close to average. February was the coldest month since February 1947 in United Kingdom and it became the 5th coldest February in the CET records dating back to 1659.
July 1986 – A polar blast in southeastern Australia brought sea level snow to the cities of Melbourne, Launceston and Hobart, with melting snow flurries even reported in parts of Sydney.
1987
January 1987 Southeast England snowfall – This was a notably cold winter month for the United Kingdom and snowy too, especially so for the southeast with a very heavy lake-effect type snow event that affected the areas of East Anglia, south-east England and London between 11 and 14 January. It was the heaviest snowfall since 1981/82.
March 1987 Greece cold wave - very rare cold air mass trapping phenomenon - The cold wave of March '87 lasted at least ten days in Athens and more in northern Greece according to weather reports and caused by a cold air mass that was trapped in the region of Greece after a high pressure system had been extended pushing trapped cold air from Russia.
1989
February 1989 featured a significant, week-long cold wave across the Western United States. Major cities affected by the cold stretched from Seattle, WA to as far south as Los Angeles, CA. Las Vegas set a record low for February with degrees on February 7, 1989.
December 1989 United States cold wave – In late 1989, the central and eastern United States saw one of the coldest Decembers on record. A white Christmas occurred.
1990–1991
Winter of 1990–91 in Western Europe – This winter was noted for its effects especially on the United Kingdom and for two significantly heavy snowfalls which occurred in December 1990 and February 1991, such snowfalls would not be seen again until February 2009. The winter was the coldest since January 1987.
December 1990 western United States – Extreme cold dropped down from Canada in the second half of December, causing record low temperatures up and down the West Coast, including one of California's most damaging freezes since 1949.
1994
1994 Northern US/Southern Canada cold outbreak – January 1994 was the coldest month ever recorded or since January 1977 or February 1934 over many parts of the northeast and north-central United States, plus adjacent southeastern Canada. Many overnight record lows were set. Cold outbreaks continued into February but the severity eased somewhat. Detroit, Michigan saw the city's coldest temperature since 1985.
1995
1995 White Earthquake in southern Chile – In August 1995 southern Chile was struck by a cold wave consisting of two successive cold fronts. Fodder scarcity caused a severe livestock starvation. Cows and sheep were also buried in snow. In parts of Tierra del Fuego up to 80% of the sheep died.
December 1995 Great Britain cold wave – On the 30th of December the United Kingdom recorded a record low of in Altnaharra in Scotland equalling the record set on February 11, 1895, and January 10, 1982.
1996
1996 Great Midwest cold outbreak – Late January and early February was Northern Minnesota's coldest short-term period on record. The record low of was recorded in Tower, Minnesota. Cities like Minneapolis experienced temperatures near .
1997
1997 Northern Plains cold air outbreak – Mid-January across the Northern U.S. was one of the windiest on record. With a low of around in some places, wind caused bitterly cold wind chills sometimes nearing . Northern parts of North Dakota saw up to of snow. This was one of the most severe cold-air outbreaks of the 1990s.
2000
July 2000 was one of the coldest months on record in South America. 4 streams of cold air pushed through the continent in a matter of 2 weeks, causing extremely low temperatures in many countries. Asunción registered -1 °C (31 °F), Buenos Aires (Ezeiza Airport) peaked at around -3 °C (27 °F), and in Brazil many cities recorded the lowest temperatures in many years. Curitiba had incredible 12 days of freezing temperatures, peaking at -4 °C (26 °F), and some areas of the southern states of Rio Grande do Sul and Santa Catarina reached -11 °C (14 °F). It also snowed heavily in these states. The cold front even reached the Amazon, with some cities in the southern part of the forest nearing temperatures of 10 °C (50 °F).
21st-century cold waves (2001–present)
2002
2002 Cold wave in Greece January 4 - the first significant snowfall of 21st century in Athens, Greece occurred on January 4, lasted three days and caused major disruption to the city. The dense snowfall during midnight of January 5 left at least 15 cm of snow on the ground in downtown Athens, the northern suburbs of Athens received more than 50 cm.
2004–2005
2004 January cold outbreak, Northeast United States – New England was close to a record month when frequent Arctic fronts caused unusually cold weather. Boston had its coldest January since 1893 (), when it averaged , and its lowest mean maximum at . Virginia Beach had an unusually long period of below freezing weather. Some areas of northern New York saw of snow in a month. Many parts of the western and midwestern area of the country seen the effect as well.
2004 February cold wave in Greece on February 12 – coldest air mass to hit Greece in 21st century till today – a very cold air mass from the North Pole was pulled in eastern Europe because of a high-pressure system in west Scandinavia extended to south Europe, resulting in a low-pressure system in the Balkans. The cold snap finally moved south hitting Greece and caused ' record low temperatures, frozen rain phenomenon, thundersnow in Athens, snowfall that blanketed Crete island, snow that fell even in the most southern place of Greece in Gavdos Island. During midnight of February 13 temperature in Athens dropped so fast that resulted in a frozen rain that ice- capped all trees, followed by a rare thundersnow phenomenon dropping 25 cm of snow in central Athens and more than 50 cm in eastern suburbs. In the morning of February 13 was recorded the lowest temperature of 21st century in the center of Athens till today of -5 °C(23 °F) degrees(some meteorological stations east of downtown on higher altitude recorded -7 °C(19,4 °F) ) resulting in water supply disruption and residents had to activate boilers and water heaters in order to make the water running again.
2004–2005 Southern Europe cold wave – All areas of Southern Europe saw an unusually hard winter. This cold front caused snow in Algeria, which is extremely unusual. The south of Spain and Morocco also recorded freezing temperatures, and record freezing temperatures were observed in the north of Portugal and Spain.
2005–2006
2005–06 European cold wave – Eastern Europe and Russia saw a very cold winter. Some of them saw their coldest on record or since the 1970s. Snow was in abundance in unusual places, such as in southern Spain and Northern Africa. All the winter months that season saw temperatures well below average across the continent.
2007
2007 Northern Hemisphere cold wave – All of Canada and most of the United States underwent a freeze after a two-week warming that took place in late March and early April. Crops froze, wind picked up, and snow drizzled much of the United States. Some parts of Europe also experienced unusual cold winter-like temperatures, during that time.
July 2007 Argentine winter storm – An interaction with an area of low pressure systems across Argentina during July 6, 7 and 8 of 2007, and the entry of a massive polar cold snap resulted in severe snowfalls and blizzards, and recorded temperatures below . The cold snap advanced from the south towards the central zone of the country, continuing its displacement towards the north during Saturday, July 7. On Monday, July 9, the simultaneous presence of very cold air, gave place to the occurrence of snowfalls. This phenomenon left at least 23 people dead.
2008
2008 Alaska cold wave – In early February, Alaska experienced the coldest temperatures for eight years, with Fairbanks nearing and Chicken, Alaska, bottoming out at , a mere away from the record of . The first half of January also brought unusual cold weather and heavy snow to widespread regions of China and the Middle East, snowfall was present in Baghdad for the first time since the 1910s.
2008 Greece cold wave February 16 - most significant cold wave that lasted three days after 2002 - On February 14 a massive high pressure system between Greenland and Scandinavia extended from Arctic Ocean to west Mediterranean accompanied by a massive low pressure system in Siberia resulting in a cold air mass in eastern Europe reaching Greece on February 16. The cold air that moved over Aegean sea caused lake-effect snow dropping more than 40 cm of snow in Athens in the Chaidari district. In addition, it caused major disruption to the city and numerous flights were cancelled due to bad weather conditions.
2009–2010
Early 2009 European Cold Wave – Early January gave most of Europe, especially in central and south very cold temperatures. Some places like Germany, France, Italy, Romania and Spain had record cold temperatures well below . Most of the places were covered in snow and ice which caused school closings and airport delays. Large cities like Paris, Madrid, Berlin and even Marseille saw very cold temperatures with much snow and ice in Northern Italy, most of Germany, in northern Portugal and even along the coasts of the Mediterranean. In early February another cold front brought heavy snowfall to much of Western Europe with the heaviest snow falling in France, Northern Italy, the Low Countries and the United Kingdom, where parts of Southern England had seen the worst snowfall in over eighteen years causing widespread travel disruption particularly around London.
February 2009 Great Britain and Ireland snowfall – The February 2009 Great Britain and Ireland snowfall was a prolonged period of snowfall that began on 1 February 2009. Some areas experienced their largest snowfall levels in 18 years. Snow fell over much of Western Europe. The United Kingdom's Met Office and Ireland's Met Éireann issued severe weather warnings in anticipation of the snowfall. More than of snow fell on parts of the North Downs and over in parts of the London area. Such snow accumulation is uncommon in London. On the morning of 6 February the majority of Great Britain and Ireland had snow cover, with the area surrounding the Bristol Channel (South Wales (Cardiff area) and South West England (Bristol area)) being most affected – had settled overnight around Okehampton, Devon, South West England with similar depths in South Wales. In Ireland the highest totals were recorded around East Kildare and Wicklow County's were up to fell around Naas, County Kildare and even more along the Wicklow Mountains. The last time such widespread snowfall affected Britain was in February 1991. On the 2nd a total of had fallen in Leatherhead, Surrey just south of the M25. Also had fallen over the South Downs and in higher areas of Brighton.
2009–10 European cold wave – At least ninety were confirmed dead after record low temperatures and heavy snowfall across Europe causes travel disruption to much of the continent including the British Isles, France, the Low Countries, Germany, Austria, Italy, Poland, the Baltic States, the Balkans, Ukraine and Russia. It was the coldest winter and longest cold spell for thirty years in the United Kingdom, whilst temperatures in the Italian Alpine peaks reached low to an extreme of .
The first snowfall began on 17 December 2009, before a respite over the Christmas period. The most severe snowy weather began on 5 January in North West England and west Scotland with temperatures hitting a low of in Greater Manchester, England. The snow spread to Southern England on 6 January and by 7 January the United Kingdom was blanketed in snow, which was captured by NASA's Terra satellite. The thaw came a week later, as temperatures started to increase. The winter weather brought widespread transport disruption, school closures, power failures, the postponement of sporting events and 25 deaths. A low of was recorded in Altnaharra, Scotland on 8 January 2010. Overall it was the coldest winter since 1978–79, with a mean temperature of .
Winter of 2009–10 in Great Britain and Ireland – The winter of 2009–10 in the United Kingdom (also called The Big Freeze by British media) was a meteorological event that started on 16 December 2009, as part of the severe winter weather in Europe. January 2010 was provisionally the coldest January since 1987 across the country. A persistent pattern of cold northerly and easterly winds brought cold, moist air to the United Kingdom with many snow showers, fronts and polar lows bringing snowy weather with it.
A cold wave affected much of the Deep South in the United States and Florida in January and February 2010.
2010–2011
Winter of 2010–11 in Great Britain and Ireland – This winter was referred to as The Big Freeze by national media in both United Kingdom and Ireland and it was the coldest winter in Britain for 31 years with an average temperature of . The UK had its second-coldest December on record, only behind December 1890, and is considered the coldest on record outside of England.
2011 New Zealand snowstorms- Caused by Antarctic storms moving upward, the whole country was affected briefly in July 2011, only for it to return even stronger in August for a prolonged period of time. Te Waipounamu/ The South Island was the hardest hit, although the generally more mild Te Ika-a-Maui/ The North Island was also affected to a large extent. Widespread and heavy snow fell in Wellington for the first time in twenty years, and in Auckland the first time since the 1930s. Despite this, the Kiwis, especially children, weren't disrupted by the snowstorms.
2012
Early 2012 European cold wave – As of February 11, 2012, at least 590 people died during a cold snap with temperatures falling below in some regions. In Ukraine, over 100 deaths were attributed to the cold.
2013
United Kingdom March–April 2013 – The UK Spring 2013 cold wave was a prolonged spell of cold weather which brought with it very heavy snowfalls, the worst in March for 30 years and since 1947 in some places. There was also some very cold temperatures with England (CET) having its coldest March since 1883 with a mean monthly temperature of . This meant that March was colder than all three winter months December 2012, January and February 2013.
Spring 2013 North American cold wave – Although the core winter of 2012–13 was fairly mild, both March and April were unusually cold across the Midwest, resulting in sharp temperature contrasts from March 2012 to March 2013 all over the United States and Canada. This late cold wave was unexpected because February and March 2013 were both forecasted to be even milder and more springlike than February and March 2012, but instead turned out with a near-average February and an unusually cold March. This same cold wave extended well into the month of April, as four notable winter storms impacted much of the northern United States, especially across Minnesota and the Dakotas. Minnesota experienced a rare May snowstorm as a result of this cold wave.
In July 2013, South America experienced the most intense cold wave in 13 years. Some coastal areas of Argentina and Uruguay had multiple days of nearly freezing temperatures, and snow fell throughout Southern Brazil, even being registered in Curitiba for the first time since 1988.
2013 Middle East cold snap - Cyprus, Egypt, Israel, West Bank, Jordan, Lebanon and Syria.
2013-2014
December 2013 North American cold wave – On December 1, the weakening of the polar vortex resulted in the jet stream shifting southward, which allowed abnormally cold temperatures to intrude the Central United States. On December 6, a daily record snowfall of was set in the Dallas–Fort Worth metroplex, breaking the old record of trace amounts of snow, set in 1950. The cold wave continued into December 10, before the temperatures returned to a more stable range.
Early 2014 North American cold wave – On January 2–11, cold arctic air initially associated with a nor'easter invaded the central and eastern United States and Canada, east of the Rockies. Temperatures were even colder than the North Pole and the South Pole in many regions in the Upper Midwest and Canada. Temperatures reached as cold as , and did not even get out of the negative double-digit temps in many places, including Chicago. The cold wave extended for a few more months, bringing a continuous pattern of record-low temperatures to most of the Central and upper eastern United States, before the pattern finally ended in early April.
2014-2015
November 2014 North American cold wave – Between November 8 and November 23, a polar vortex similar to earlier in 2014 has a temporary comeback, delivering the 2014–15 winter season's first three significant winter storms in the United States. Snowfall records were confirmed all over the Midwest and the Northeast, especially around the Great Lakes. Buffalo, New York, was among the hardest hit in the unseasonably wintry November.
February 2015 North American cold wave – During the second half of February 2015, temperature records were broken in both sides of the spectrum. Extreme warm records were broken in the western half of the United States and extreme cold records were broken in the eastern half. In addition to the extreme cold wave at its most brutal in the Great Lakes, Mid-Atlantic, and New England, snowfall was reported as far south as Tupelo, Mississippi; Huntsville, Alabama; and Shreveport, Louisiana. The cold wave became widespread and all the remaining mild conditions from the west were pushed into northern Mexico. The cold wave even extended well into early March, with a part of every U.S. state except Florida reporting a snow cover by March 1, 2015.
2016
January 2016 East Asia cold wave – Caused over 100 known deaths across East Asia, South Asia and Southeast Asia.
February 2016 North American cold snap – A cold wave hit North America during the second week of February 2016 which caused record-breaking wind chills and temperatures. New York City broke its record low of -1 °F (−18 °C), which was the first sub-zero reading for New York since January 19, 1994. On 13 February 2016, Whiteface Mountain underwent a record windchill of −114 °F (−81 °C), while in Boston, Massachusetts, the temperature dropped to −9 °F (−23 °C), the coldest since 1957.
2017
April 2017 - During the same year, low temperatures were recorded in Sarajevo from an unusual snowfall that hit the city disrupting the traffic since the 2012 cold wave.
January 2017 European cold wave – A cold wave hit Central and East Europe on January 5. The lowest temperature was degrees. The cold caused at least 60 deaths. There was also massive snowfall.
2017–2018
Cold wave of November 9–12, 2017. Record lows were broken from Minneapolis to Washington, D.C., as Arctic air swept through the areas.
Cold wave starting late December 2017 (December 24 respectively), North America. A persistent wave of temperature extremes, including a cold wave, took place in Canada and the northeastern and central areas of the United States from Northern Canada to Mississippi, with temperatures in much of Canada of around and as low as in New York state, and as high as and in Sandberg and Los Angeles, respectively, in California.
2018
Late February and early–mid March 2018, Europe. Easterly winds created bitter conditions, along with snow, twice during early spring 2018. Storm Emma, which affected southern areas of Great Britain, brought up to of snow. Overall, 93 people across Europe died; 27 deaths occurred in Poland and 17 in the UK.
Around early December 2018, some cold waves hit Central Russia and Kazakhstan, bringing temperatures several degrees down from the average. Big cities like Novosibirsk, Krasnoyarsk, Omsk, Irkutsk, and Barnaul had experienced temperatures of -10 C or lower. The cold wave dissipated in mid-December, but returned close to Christmas, impacting Central Russia and Kazakhstan strongly again, before finally dissipating around New Year's Eve. The lowest temperature, happened in Yerbogachen on Christmas.
2019
In late January, an extreme cold wave hit Canada and the midwest of the United States, bringing temperatures below , with all-time record lows set in several cities. New York City had a low temperature of 2 °F (-16 °C), the coldest reading in Central Park since February 14, 2016, when the mercury dropped to -1 °F (-18 °C), the coldest sub-zero reading in New York City since January 19, 1994.
February 2019: Brought temperatures 10~15 degrees lower than the February average lows, and temperatures of degrees to many parts of Siberia again. Novosibirsk, the largest city in Asian Russia, lowered to on 2 February, just a fraction of a degree shy to its previous record in 1977. Krasnoyarsk also lowered to on 4 February, missing just 0.3 °C to its record in 2001. Irkutsk Oblast had also recorded very low temperatures, with Irkutsk hitting on 6 February, beating its previous record of 6 February, Bratsk has seen temperatures below on February 5, which is sometimes colder than cities like Yakutsk. 7 and 8 February was even more brutal, with Omsk as low as , and a merciless record low in Novokuznetsk, just 0.3 °C short of the record low in 1969. Lowest temperature in the cold wave of was recorded in Vanavara The cold wave also lightly affected the Russian Far East and some parts North America. California experienced an unusually wet and cold February, where Los Angeles experienced its coldest February since 1962.
November 2019: North American cold wave
November 2019: Turkish cold wave
2020
February 2020: 2nd time snow falls in Baghdad since the 1910s.
October 2020: The Rocky Mountains and Great Plains of the United States experienced record-cold temperatures. Subzero Fahrenheit temperatures (< -18 °C) were recorded in Montana, Wyoming, Colorado, and Nebraska. Bozeman, Montana, reached and Denver, Colorado reached .
2021
January 2021: Severe cold wave hit many regions in Eurasia, especially in the Iberian Peninsula, Central Asia, South Asia, East Asia and Southeast Asia with some of the lowest temperatures in many years. In Spain, Clot del Tuc de la Llança in Catalonia in the Spanish Pyrenees recorded −34.1 °C on January 6. It was not only the new national lowest temperature record in Spain, 2 °C lower than the previous record in 1956, but also the lowest temperature in the whole Iberian Peninsula ever recorded. On 9 January another new record was attained with -35.6 °C in Vega de Liordes, in the Cantabrian mountains of Spain. The cold wave was concurrent with a historic snowstorm which covered Madrid under 50 to 60 cm of snow, the first time since 1971. Novosibirsk reached a low of -41 °C. Beijing recorded a low of -19.6 °C which was the coldest since 1966. Seoul also recorded -18.6 °C in January 8 which was the tie record with 2001 and the coldest since 1986. Over 200 cm of snow fell in western Japan along the Japan Sea coast.
February 2021: Turkish cold wave
February 2021: February 2021 North American cold wave
May 2021: Sydney CBD recorded its coldest stretch of May days in 54 years due to a polar blast that swept across Australia's southeast, which kept the temperatures below in the early mornings for five consecutive nights, in addition to the inland suburbs of Sydney dipping down to in Camden and in Sydney Olympic Park.
June 2021, Sydney CBD had its coldest day since 1984 and the coldest June day since 1899, where it reached a maximum of just . Bankstown, a western suburb, only reached , its coldest day in 50 years, with nearby suburbs registering a similar temperature. These unusually cold maximums were caused by a cut-off low.
2022
February-March 2022: Turkish cold wave
May 2022: 2022 South American cold wave
June 2022: Aysén Region in southern Chile experienced unusually cold temperatures. In numerous locations, temperatures dropped below -10 C. In one case it reached -14 C. Parts of Tortel Fjord and Simpson River froze. In Los Lagos Region, north of Aysén Region, temperatures of -5 C were registered.
December 2022: Late December 2022 North American winter storm
December 2022: December 2022 United Kingdom cold wave - A relatively short-lived but significant cold spell that brought some of the lowest temperatures recorded in the United Kingdom since the winter 2010-11.
2023
February 2023: February 2023 North American cold wave
2024
January 2024: Mid-January 2024 North American Cold Wave - Freezing temperatures affected campaign events leading up to the Iowa Caucuses on January 15, and could affect turnout in the first contest in the 2024 presidential primaries and caucuses. In addition, the cold temperatures resulted in the 4th coldest NFL game on record between the Kansas City Chiefs and Miami Dolphins, with kickoff temperature being . 15 people were hospitalized due to the cold temperatures during the game. Wind chills following the storm reached as low as in Montana, and wind chills were still as far south as Dallas on January 14. On January 13, Dillon, Montana reached an all time record low of , while Bozeman, Montana recorded their second coldest temperature at . That same day, Dickinson, North Dakota reached a wind chill, their coldest since the wind chill formula was updated in 2001, and a air temperature, a daily record and their coldest temperature since 1990. On January 16, Houston dropped to a daily record low of .
| Physical sciences | Seasons | Earth science |
19625172 | https://en.wikipedia.org/wiki/Ocean%20surface%20topography | Ocean surface topography | Ocean surface topography or sea surface topography, also called ocean dynamic topography, are highs and lows on the ocean surface, similar to the hills and valleys of Earth's land surface depicted on a topographic map.
These variations are expressed in terms of average sea surface height (SSH) relative to Earth's geoid. The main purpose of measuring ocean surface topography is to understand the large-scale ocean circulation.
Time variations
Unaveraged or instantaneous sea surface height (SSH) is most obviously affected by the tidal forces of the Moon and by the seasonal cycle of the Sun acting on Earth. Over timescales longer than a year, the patterns in SSH can be influenced by ocean circulation. Typically, SSH anomalies resulting from these forces differ from the mean by less than ± at the global scale. Other influences include changing interannual patterns of temperature, salinity, waves, tides and winds. Ocean surface topography can be measured with high accuracy and precision at regional to global scale by satellite altimetry (e.g. TOPEX/Poseidon).
Slower and larger variations are due to changes in Earth's gravitational field (geoid) due to melting ice, rearrangement of continents, formation of sea mounts and other redistribution of rock. The combination of satellite gravimetry (e.g. GRACE and GRACE-FO) with altimetry can be used to determine sea level rise and properties such as ocean heat content.
Applications
Ocean surface topography is used to map ocean currents, which move around the ocean's "hills" and "valleys" in predictable ways. A clockwise sense of rotation is found around "hills" in the northern hemisphere and "valleys" in the southern hemisphere. This is because of the Coriolis effect. Conversely, a counterclockwise sense of rotation is found around "valleys" in the northern hemisphere and "hills" in the southern hemisphere.
Ocean surface topography is also used to understand how the ocean moves heat around the globe, a critical component of Earth's climate, and for monitoring changes in global sea level.
The collection of the data is useful for the long-term information about the ocean and its currents. According to NASA science this data can also be used to provide understanding of weather, climate, navigation, fisheries management, and offshore operations. Observations made about the data are used to study the oceans tides, circulation, and the amount of heat the ocean contains. These observations can help predict short and long term effects of the weather and the earth's climate over time.
Measurement
The sea surface height (SSH) is calculated through altimetry satellites using as a reference surface the ellipsoid, which determine the distance from the satellite to a target surface by measuring the satellite-to-surface round-trip time of a radar pulse. The satellites then measure the distance between their orbit altitude and the surface of the water. Due to the differing depths of the ocean, an approximation is made. This enables data to be taken precisely due to the uniform surface level. The satellite's altitude then has to be calculated with respect to the reference ellipsoid. It is calculated using the orbital parameters of the satellite and various positioning instruments. However, the ellipsoid is not an equipotential surface of the Earth's gravity field, so the measurements must be referenced to a surface that represents the water flow, in this case the geoid. The transformations between geometric heights (ellipsoid) and orthometric heights (geoid) are performed from a geoidal model. The sea surface height is then the difference between the satellite's altitude relative to the reference ellipsoid and the altimeter range. The satellite sends microwave pulses to the ocean surface. The travel time of the pulses ascending to the oceans surface and back provides data of the sea surface height.
In the image below you can see the measurement system using by the satellite Jason-1.
Satellite missions
Currently there are nine different satellites calculating the earth ocean topography, Cryosat-2, SARAL, Jason-3, Sentinel-3A and Sentinel-3B, CFOSat, HY-2B and HY-2C, and Sentinel-6 Michael Freilich (also called Jason-CS A). Jason-3 and Sentinel-6 Michael Freilich are currently both in space orbiting Earth in a tandem rotation. They are approximately 330 kilometers apart.
Ocean surface topography can be derived from ship-going measurements of temperature and salinity at depth. However, since 1992, a series of satellite altimetry missions, beginning with TOPEX/Poseidon and continued with Jason-1, Ocean Surface Topography Mission on the Jason-2 satellite, Jason-3 and now Sentinel-6 Michael Freilich have measured sea surface height directly. By combining these measurements with gravity measurements from NASA's Grace and ESA's GOCE missions, scientists can determine sea surface topography to within a few centimeters.
Jason-1 was launched by a Boeing Delta II rocket in California in 2001 and continued measurements initially collected by TOPEX/Poseidon satellite, which orbited from 1992 up until 2006. NASA and CNES, the French space agency, are joint partners in this mission.
The main objectives of the Jason satellites is to collect data on the average ocean circulation around the globe in order to better understand its interaction with the time varying components and the involved mechanisms for initializing ocean models. To monitor the low frequency ocean variability and observe the season cycles and inter-annual variations like El Niño and La Niña, the North Atlantic oscillation, the pacific decadal oscillation, and planetary waves crossing the oceans over a period of months, then they will be modeled over a long period of time due to the precise altimetric observations. It aims to contribute to observations of the mesoscale ocean variability, affecting the whole oceans. This activity is especially intense near western boundary currents. Also monitor the average sea level because it is a large indicator of global warming through the sea level data. Improvement of tide modeling by observing more long period components such as coastal interactions, internal waves, and tidal energy dissipation. Finally the satellite data will supply knowledge to support other types of marine meteorology which is the scientific study of the atmosphere.
Jason-2 was launched on June 20, 2008, by a Delta-2 rocket out of the California site in Vandenberg and terminated its mission on October 10, 2019. Jason-3 was launched on January 16, 2016 by a Falcon-9 SpaceX rocket from Vandenberg, as well as Sentinel-6 Michael Freilich, launched on November 21, 2020.
The long-term objectives of the Jason satellite series are to provide global descriptions of the seasonal and yearly changes of the circulation and heat storage in the ocean. This includes the study of short-term climatic changes such as El Nino, La Nina. The satellites detect global sea level mean and record the fluctuations. Also detecting the slow change of upper ocean circulation on decadal time scales, every ten years. Studying the transportation of heat and carbon in the ocean and examining the main components that fuel deep water tides. The satellites data collection also helps improve wind speed and height measurements in current time and for long-term studies. Lastly improving our knowledge about the marine geoid. The first seven months Jason-2 was put into use it was flown in extreme close proximity to Jason-1. Only being one minute apart from each other the satellites observed the same area of the ocean. The reason for the close proximity in observation was for cross-calibration. This was meant to calculate any bias in the two altimeters. This multiple month observation proved that there was no bias in the data and both collections of data were consistent.
A new satellite mission called the Surface Water Ocean Topography Mission has been proposed to make the first global survey of the topography of all of Earth's surface water—the ocean, lakes and rivers. This study is aimed to provide a comprehensive view of Earth's freshwater bodies from space and more much detailed measurements of the ocean surface than ever before.
| Physical sciences | Oceanography | Earth science |
19629283 | https://en.wikipedia.org/wiki/Human%20anus | Human anus | In humans, the anus (: anuses or ani; from Latin ānus, "ring", "circle") is the external opening of the rectum located inside the intergluteal cleft. Two sphincters control the exit of feces from the body during an act of defecation, which is the primary function of the anus. These are the internal anal sphincter and the external anal sphincter, which are circular muscles that normally maintain constriction of the orifice and which relax as required by normal physiological functioning. The inner sphincter is involuntary and the outer is voluntary. Above the anus is the perineum, which is also located beneath the vulva or scrotum.
In part owing to its exposure to feces, a number of medical conditions may affect the anus, such as hemorrhoids. The anus is the site of potential infections and other conditions, including cancer (see anal cancer).
With anal sex, the anus can play a role in sexuality. Attitudes toward anal sex vary, and it is illegal in some countries. The anus is often considered a taboo part of the body, and is known by many, usually vulgar, slang terms. Some sexually transmitted infections including HIV/AIDS and anal warts can be spread via anal sex.
Structure
The anus is the final part of the gastrointestinal tract, and directly continues from the rectum, passing through the pelvic floor. The top and bottom of the anus are surrounded by the internal and external anal sphincters, two muscular rings which control defecation.
The anus is surrounded in its length by folds called anal valves, which converge at a line known as the pectinate line. This represents the point of transition between the hindgut and the ectoderm in the embryo. Below this point, the mucosa of the internal anus becomes skin. The pectinate line is also the division between the internal and external anus.
The anus receives blood from the inferior rectal artery and innervation from the inferior rectal nerves, which branch from the pudendal nerve.
Microanatomy
The pseudostratified columnar epithelium of the gastrointestinal tract transitions to stratified squamous epithelium at the pectinate line. The stratified squamous epithelium gradually accumulates sebaceous and apocrine glands.
Development
During puberty, as testosterone triggers androgenic hair growth on the body, pubic hair begins to appear around the anus. Although initially sparse, it fills out by the end of puberty, if not earlier. In some genetic populations, androgenic hair is less common.
Function
Defecation
Intra-rectal pressure builds as the rectum fills with feces, pushing the feces against the walls of the anal canal. Contractions of abdominal and pelvic floor muscles can create intra-abdominal pressure, which further increases intra-rectal pressure. The internal anal sphincter (an involuntary muscle) responds to the pressure by relaxing, thus allowing the feces to enter the canal. The rectum shortens as feces are pushed into the anal canal and peristaltic waves push the feces out of the rectum. Relaxation of the internal and external anal sphincters allows the feces to exit from the anus, finally, as the levator ani muscles pull the anus up over the exiting feces.
Clinical significance
Anal fissures are tears in the external lining of the lining (mucosa) of the anus. These are exquisitely painful, with pain occurring after a motion is passed; other symptoms may include minor bleeding, discharge, or itch. Generally, fissures are due to injury to the mucosa, or because of a poor local blood supply that prevents proper healing, with spasm of the external anal sphincter contributing. The external anal sphincter can be relaxed by the application of glyceryl trinitrate creams, and constipation is managed with laxatives and improving hydration. Some fissures may require botulinum toxin injection; worst cases may require surgical intervention, such as "lateral internal anal sphincterotomy or advancement anoplasty".
Hemorrhoids are visible blood vessels from the internal or external venous plexuses of the anus. Haemorrhoids may; cause bleeding after passing a motion; be painful; cause an itch; or out of the anus. Haemorrhoids are often associated with straining due to constipation, and pregnancy. Usually, haemorrhoids are managed with medications to make motions more soft and prevent straining during constipation. Some haemorrhoids require surgery to manage, which may involve placing a band around the haemorrhoid, in order for it to lose blood supply; or surgical excision.
Other
Fistula
Birth defects, including imperforation, stenosis, Tailgut cyst
Infections
Anal abscesses usually result from infection of the normal glands of the anus, or sometimes because of Crohn's disease. They usually occur to the sides of the sphincters, and between the internal and external sphincters, either on the surface, or deeper. They may get bigger, enlarging in the direction of the rectum, and resulting in an abnormal connection called an anorectal fistula. They are usually managed with surgical drainage and antibiotics.
Additional
Sexually transmitted infections
Anal warts, also called anal condyloma
Cancer
Anal cancer, also called "anal carcinoma", and anal intraepithelial neoplasia.
Itching, incontinence and constipation
Itchiness, called pruritus ani, can affect the anus area. It is most often due to long-term exposure of the anus to faeces, with reasons including diseases of the anus such as haemorrhoids, fistulas and fissures; poor hygiene or chronic diarrhoea; local infections such as tapeworm and thrush; or skin conditions such as psoriasis and contact dermatitis. If there is a specific cause identified, the cause may be treated to relieve the itch. Otherwise, treatment includes keeping the area clean and dry, ceasing topical creams and ointments, and potentially bulk-forming laxatives to reduce the chance of faecal contamination.
Damage or injury to the anal sphincter (patulous anus in more severe cases) as a result of damage during surgery, such as to the perineal region, or resulting from anal sex; can lead to flatus and/or fecal incontinence, chronic constipation and megacolon.
Society and culture
Sexuality
The anus has a relatively high concentration of nerve endings and can be an erogenous zone, which can make anal intercourse pleasurable if performed properly. The pudendal nerve that branches to supply the external anal sphincter also branches to the dorsal nerve of the clitoris and the dorsal nerve of the penis.
In addition to nerve endings, pleasure from anal intercourse may be aided by the close proximity between the anus and the prostate for males, and vagina, clitoral legs and anal area for females. This is because of indirect stimulation of the prostate and vagina or clitoral legs. For a male insertive partner, the tightness of the anus can be a source of pleasure via the tactile pressure on the penis. Pleasure from the anus can also be achieved through anal masturbation, anal fingering, facesitting, anilingus, and other penetrative and non-penetrative acts. Anal stretching or fisting is pleasurable for some, but it poses a more serious threat of damage due to the deliberate stretching of the anal and rectal tissues; its injuries include anal sphincter lacerations and rectal and sigmoid colon (rectosigmoid) perforation, which might result in death. Lubricant and condoms are widely regarded as a necessity while performing anal sex, as well as a slow and cautious penetration.
Anal intercourse is sometimes referred to as sodomy or buggery, and is considered taboo in a number of legal systems. It has been, and in some jurisdictions continues to be, a crime carrying severe punishment.
Hygiene
To prevent diseases of the anus and to promote general hygiene, humans often clean the exterior of the anus after emptying the bowels. A rinse with water from a bidet or a wipe with toilet paper is often used for this purpose, though anal cleansing practices vary greatly between cultures.
Cosmetics
Shaving, trimming, depilatory (hair removal), or Brazilian waxing can clear the perineum of hair.
Anal bleaching is a process in which the anus and perineum is lightened. Perineum sunning is a process in which the anus is sun tanned by deliberate exposure to sunlight, resulting in a darkening of the area.
A true anal piercing is rare because it may interfere with the function of the anus and cause infections. Surface piercings of the perineum are easier to care for and much more common.
Some people have their anuses tattooed.
Slang
The anus has many slang terms including asshole, butthole (and their respective British equivalents arsehole, bumhole), cornhole, and bunghole.
Additional images
| Biology and health sciences | Human anatomy | Health |
19629444 | https://en.wikipedia.org/wiki/Kagu | Kagu | The kagu or cagou (Rhynochetos jubatus) is a crested, long-legged, and bluish-grey bird endemic to the dense mountain forests of New Caledonia. It is the only surviving member of the genus Rhynochetos and the family Rhynochetidae, although a second species has been described from the fossil record.
Measuring in length, it has pale grey plumage and bright red legs. Its 'nasal corns' are a unique feature not shared with any other bird. Almost flightless, it spends its time on or near the ground, where it hunts its invertebrate prey, and builds a nest of sticks on the forest floor. Both parents share incubation of a single egg, as well as rearing the chick. It has proven vulnerable to introduced predators and is threatened with extinction.
Etymology
The name kagu is derived from the Melanesian names for the species. The species is variously known as the kavu or kagou in the Kanak languages, and as the cagou in French (also used as an alternative spelling in English).
Taxonomy and systematics
The generic name Rhynochetos, and the clade name Rhynochetidae, are derived from the Greek rhis meaning "nose" and khaitē meaning "long hair", referring to its long stiff hairs over the nostrils. The specific name jubatus is derived from the Latin iubātus meaning crested.
The kagu's affinities are not well resolved. It was long one of the most enigmatic birds and in more recent times is usually affiliated with the Gruiformes. It was initially classed as a member of the clade Ardeae because of the presence of powder down, similarities in plumage colour and internal anatomy, the colour of the chicks and eggs, and the change in colouration of the chick as it grows.
When seen as a gruiform, the kagu is generally considered related to the extinct adzebills from New Zealand and the sunbittern from Central and South America. Recent studies do indicate that the sunbittern is the closest living relative of the kagu. For example, Fain & Houde found these to be certainly sister taxa. They and the mesites did not group with traditional Gruiformes in their study, but instead with their proposed clade Metaves, which also includes the hoatzin, pigeons, nightjars, flamingos, tropicbirds, Apodiformes, sandgrouse, and grebes. The internal structure of this group was not well resolvable by their data, although later studies confirmed a close relationship between the kagu and sunbittern.
The kagu and sunbittern, and possibly the adzebills, seem to form a distinct Gondwanan lineage of birds, either as one order or possibly more. Although the relationships between them and groups previously considered related, such as the mesites and the "core Gruiformes," are not yet resolved. It is notable, however, that the sunbittern and the mesites possess powder down too, whereas the "core Gruiformes" do not. The ancestors of the kagu are believed to have diverged from the sunbittern in the Oligocene, 45 to 17 million years ago, and colonized New Caledonia 60 to 25 million years ago. In the absence of terrestrial predators, it eventually became flightless.
While the kagu is the only living species in the clade Rhynochetidae, a larger species, the lowland kagu (Rhynochetos orarius), has been described from Holocene subfossil remains. The measurements of this species were 15% bigger than Rhynochetos jubatus, with no overlap in measurements except those of the forelimbs. Given that the sites from which R. orarius remains have been recovered are all lowland sites, and that no fossils of R. jubatus have been found in these sites, the scientists that described the fossils have suggested that they represent highland and lowland species respectively. R. orarius is one of many species to have become extinct in New Caledonia after the arrival of humans. The validity of the species has been questioned by some authors, but accepted by others.
Description
The kagu is a ground-living bird, in length. The weight can vary considerably by individual and by season, ranging from . Its plumage is unusually bright for a bird of the forest floor; ash-grey and white coloured. There is little sexual dimorphism beyond a difference in the amount of barring in the primary feathers. It possesses powder down which helps keep it dry and insulates it in the extremes of New Caledonia's tropical climate. The crest, which is used to display to other members of the species, is barely noticeable when at rest but can be erected and fanned out.
It is nearly flightless, using its wings for displays (its primary wing feathers are patterned), and for moving quickly through the forest. It can also use them to glide when fleeing danger. The wings are not reduced in size like some other flightless birds, and have a span of around , but they lack the musculature for flight. These wings are also used for a 'broken-wing' display, a behaviour shared with their relative sunbittern, used to fake an injury and draw the attention of a predator away from their chick. It possesses bright red legs which are long and strong, enabling the bird to travel long distances on foot and run quickly.
It has large eyes, positioned so that they give good binocular vision which is helpful in finding prey in the leaf litter and seeing in the gloom of the forest. It possesses 'nasal corns', structures covering its nostrils, which are a feature not shared by any other bird. These are presumed to prevent particles entering the nostrils when probing in soil during feeding. Another unique characteristic of the species is that it has only one-third as many red blood cells and three times more hemoglobin per red blood cell than is usual in birds.
Distribution and habitat
The kagu is endemic to the forests and shrubland of New Caledonia. Within that island group it is restricted to the main island of Grande Terre. There is no evidence that it occurred on the Loyalty Islands, although fossil remains of the extinct lowland form R. orarius have been found on the Ile des Pines. The kagu is a habitat generalist and able to exist in a range of different forest types if sufficient prey is present, from rain forest to drier lowland forest. They are also able to feed in some drier shrubland associated with the island's ultrabasic rocks, although not in the poorest, low-prey shrubland of this type. They are also absent from areas where extensive ground cover makes foraging difficult, such as grassland or areas with high fern cover, but may pass through such areas to reach other foraging areas. The species has undergone some range contraction due to hunting and predation by introduced species. Its original, pre-human distribution, and the extent to which it and its sister species R. orarius coexisted in lowland areas of New Caledonia, is still not fully understood and awaits further research into the subfossil record.
Behaviour and ecology
Kagu are territorial, maintaining year-round territories of around . They have a clan-based social organization, with families composed of one breeding female and one to three breeding males. Male offspring also help to defend their parents' territorial claims. However the social organisation of kagu has been disrupted in recent years due to attacks by dogs and families. Cases where either the breeding male or female have been killed have led to non-fraternal polyandrous behaviour. Cooperative and unrelated polyandry is rare in birds, but has been seen in species such as the dunnock and the Tasmanian nativehen.
Within the territory the pairs are solitary during the non-breeding season, and may have separate but overlapping foraging areas. Kagus make a range of different sounds, most commonly while duetting in the morning, each duet lasting about 15 minutes. The kagu's crest and wings are used in territorial displays towards other kagu, slightly different displays are used towards potential predators. Territorial disputes may be resolved by fighting using wings and bills; in the wild this seldom results in serious injuries.
Diet
The kagu is exclusively carnivorous, feeding on a variety of animals, with annelid worms, snails and lizards being amongst the most important prey items. Also taken are larvae, spiders, centipedes and insects, such as grasshoppers, bugs, and beetles. The majority of the diet is obtained from the leaf litter or soil, with other prey items found in vegetation, old logs and rocks. Sometimes kagus will hunt small animals in shallow water. Their hunting technique is to stand motionless on the ground or from an elevated perch, and silently watch for moving prey. They may stand on one foot and gently move the leaf litter with the other foot in order to flush prey. Having located prey they will move towards the prey and stand over it, ready to strike, or make a dash towards the prey from their watching location. If digging is required to obtain the prey this is done with the bill, the feet are not used to dig or scratch away debris.
Breeding
Kagus are monogamous breeders, generally forming long-term pair bonds that are maintained for many years, even possibly life. Kagu can be long lived, with birds in captivity living for over 20 years. A single nesting attempt is made each year, although should the first nesting attempt fail a second attempt is made that year. A simple nest is constructed, which is little more than a heaped pile of leaves, although in some cases the egg may be laid directly on the ground. The nest is not concealed but is usually adjacent to a tree trunk, log or low vegetation. A single grey slightly blotched egg is laid which weighs 60–75 g. Incubation duties are shared by the parents. Each bird will incubate the egg for 24 hours, with the changeover occurring around noon each day. During each incubation stint the parent will remain on the egg the whole time except early in the morning, when the bird will briefly move away to call to its mate and occasionally forage quickly. The incubation period lasts for 33–37 days, which is long for the size of the egg.
Offspring may remain in their parents' territory for many years after fledging, sometimes up to six years. These chicks do not help in incubating the eggs or raising the chicks, but nevertheless improve the breeding success of the parents. The older offspring do apparently help in territory defence, responding to playback of rivals and also participating in territorial fights, and it has been suggested that this should be treated as a form of cooperative breeding. The kagu has also been observed adopting an unrelated chick, a behaviour more common in species with low reproductive output, high social organisation, and extensive parental care of the young, all traits shared by the kagu.
Status and conservation
The kagu's initial decline was caused by subsistence hunting. The bird was trapped extensively for the European pet trade and for museums and zoos until it was afforded protection. It is threatened by introduced cats, pigs and dogs. New Caledonia lacked mammals (except for bats) before the arrival of humans, and many of its native species have been negatively affected by introduced mammals. Rats have a big impact on nestlings, accounting for 55% of nestling losses. Kagus also suffer from habitat loss caused by mining and forestry.
Concern was first raised about the future of the kagu in 1904. A visiting American scientist noted in 1948 that the extinction of the species was probable, and identified the many threats the species faced. The first concrete evidence of the impact of dogs came when a New Zealand researcher's study population was quickly exterminated by dogs in the 1990s, although suspicions about the importance of dogs and other predators had been voiced before this and dog control measures had been enacted in some areas in the 1980s. The kagu is listed as endangered (CITES I) and enjoys full protection in New Caledonia. It has been the subject of dedicated conservation efforts and is receptive to ex-situ conservation, breeding well in Nouméa Zoo. It is also prospering in Rivière Bleue Territorial Park, which has a pest-management programme and has been the site of releases into the wild of captive-bred birds.
Recent research has shown that naturally occurring heavy metals in the soil may affect Kagu through their food supply. Kagu in areas where soil levels of heavy metals were low laid more eggs and had higher numbers of fledglings, as well as having smaller home-ranges and higher body mass, than Kagu in areas where the soil was heavy-metal rich. It has therefore been suggested that Kagu conservation is likely to be more effective in areas where heavy-metal levels in the soil are low.
Relationship with humans
The kagu had an important role in the traditional lives of the Kanak tribes of New Caledonia. Among the tribes found in the vicinity of Hienghène in the north of Grande Terre, its name was given to people, its crest was used in the head-dresses of chiefs, and its calls were incorporated into war dances and considered messages to be interpreted by the chiefs. Kanaks in the vicinity of Houaïlou referred to the species as the "ghost of the forest."
The species was not discovered by Europeans until the French colonisation of New Caledonia in 1852 and was not described until a specimen was taken to the Colonial Exhibition in Paris in 1860. This led to a surge in scientific interest in the species, which resulted in many birds being trapped for museums and zoos. The species was also trapped for food and was considered a delicacy by European colonisers. It was also fashionable to own kagus as pets. A campaign was run from 1977–1982 to phase out the pet trade in kagus. Today, the kagu is considered very important in New Caledonia; it is a high-profile endemic emblem for the territory. Its distinctive song used to be played to the nation every night as the island's TV station signed off the air. Its survival is considered important for the territory's economy and image.
| Biology and health sciences | Basics | Animals |
19630671 | https://en.wikipedia.org/wiki/Laurentia | Laurentia | Laurentia or the North American Craton is a large continental craton that forms the ancient geological core of North America. Many times in its past, Laurentia has been a separate continent, as it is now in the form of North America, although originally it also included the cratonic areas of Greenland and the Hebridean Terrane in northwest Scotland. During other times in its past, Laurentia has been part of larger continents and supercontinents and consists of many smaller terranes assembled on a network of early Proterozoic orogenic belts. Small microcontinents and oceanic islands collided with and sutured onto the ever-growing Laurentia, and together formed the stable Precambrian craton seen today.
Etymology
The craton is named after the Laurentian Shield, through the Laurentian Mountains, which received their name from the St. Lawrence River, named after Saint Lawrence of Rome.
Interior platform
In eastern and central Canada, much of the stable craton is exposed at the surface as the Canadian Shield, an area of Precambrian rock covering over a million square miles. This includes some of the oldest rock on Earth, such as the Archean rock of the Acasta Gneiss, which is 4.04 billion years (Ga) old, and the Istaq Gneiss Complex of Greenland, which is 3.8 Ga. When subsurface extensions are considered, the wider term Laurentian Shield is more common, not least because large parts of the structure extend outside Canada. In the United States, the craton bedrock is covered with sedimentary rocks on the broad interior platform in the Midwest and Great Plains regions and is exposed only in northern Minnesota, Wisconsin, the New York Adirondacks, and the Upper Peninsula of Michigan. The sequence of sedimentary rocks varies from about 1,000 m to in excess of 6,100 m (3,500–20,000 ft) in thickness. The cratonic rocks are metamorphic or igneous with the overlying sedimentary layers composed mostly of limestones, sandstones, and shales. These sedimentary rocks were largely deposited 650–290 Ma.
The oldest bedrock, assigned to the Archean Slave, Rae, Hearne, Wyoming, Superior, and Nain Provinces, is located in the northern two thirds of Laurentia. During the Early Proterozoic they were covered by sediments, most of which has now been eroded away.
Greenland is part of Laurentia. The island is separated from North America by the Nares Strait, but this is a Pleistocene erosional feature. The strait is floored with continental crust and shows no indications of a thermal event or seaway tectonism. Greenland is composed mostly of crust of Archean to Proterozoic age, with lower Paleocene shelf formations on its northern margin and Devonian to Paleogene formations on its western and eastern margins. The eastern and northern margins were heavily deformed during the Caledonian orogeny.
The Isua Greenstone Belt of western Greenland preserves oceanic crust containing sheeted dike complexes. These provide evidence to geologists that mid-ocean ridges existed 3.8 Ga. The Abitibi gold belt in the Superior Province is the largest greenstone belt in the Canadian Shield.
Tectonic history
Assembly
Laurentia first assembled from six or seven large fragments of Archean crust at around 2.0 to 1.8 Gya. The assembly began when the Slave craton collided with the Rae-Hearne craton, and the Rae-Hearne craton collided shortly after with the Superior Craton. These then merged with several smaller fragments of Archean crust, including the Wyoming, Medicine Hat, Sask, Marshfield, and Nain blocks. This series of collisions raised the mountains of the Trans-Hudson orogenic belt, which likely were similar to the modern Himalayas, and the Wopmay orogen of northwest Canada. During the assembly of the core of Laurentia, banded iron formation was deposited in Michigan, Minnesota, and Labrador.
The resulting nucleus of Laurentia was mostly reworked Archean crust but with some juvenile crust in the form of volcanic arc belts. Juvenile crust is crust formed from magma freshly extracted from the Earth's mantle rather than recycled from older crustal rock. The intense mountain building of the Trans-Hudson orogeny formed thick, stable roots beneath the craton, possibly by a process of "kneading" that allowed low density material to move up and high density material to move down.
Over the next 900 million years, Laurentia grew by the accretion of island arcs and other juvenile crust and occasional fragments of older crust (such as the Mojave block). This accretion occurred along the southeastern margin of Laurentia, where there was a long-lived convergent plate boundary. Major accretion episodes included the Yavapai orogeny at 1.71 to 1.68 Gya, which welded the 1.8 to 1.7 Gya Yavapai province to Laurentia; the Mazatzal orogeny at 1.65 to 1.60 Gya, accreting the 1.71 to 1.65 Gya Mazatzal province; the Picuris orogeny at 1.49 to 1.45 Gya, which may have welded the 1.50 to 1.30 Gya Granite-Rhyolite province to Laurentia; and the Grenville orogeny at 1.30 to 0.95 Gya, which accreted the 1.30 to 1.00 Gya Llano-Grenville province to Laurentia.
The Picuris orogeny, in particular, was characterized by the intrusion of great volumes of granitoid magma into the juvenile crust, which helped mature the crust and stitch it together. Slab rollback at 1.70 and 1.65 Gya deposited characteristic quartzite-rhyolite beds on the southern margin of the craton. This long episode of accretion doubled the size of Laurentia but produced craton underlain by relatively weak, hydrous, and fertile (ripe for extraction of magma) mantle lithosphere. The subduction under the southeast margin of the continent likely caused enrichment of the lithospheric mantle beneath the orogenic belts of the Grenville Province. Around 1.1 Gya, the center of the craton nearly rifted apart along the Midcontinent Rift System. This produced the Keweenawan Supergroup, whose flood basalts are rich in copper ore.
Formation and breakup of Rodinia
Laurentia was formed in a tectonically active world. The subduction under the southeast margin of the continent is thought to have contributed to the formation of Rodinia. According to the Southwest U.S. and East Antarctica or SWEAT hypothesis, Laurentia became the core of the supercontinent. It was rotated approximately 90 degrees clockwise compared with its modern orientation, with East Antarctica and Australia to the north (what is now the west), Siberia to the east (present north), Baltica and Amazonia to the south (present east), and Congo to the southwest (present southeast). The Grenville orogen extended along the entire southwest (present southeast) margin of Laurentia, where it had collided with Congo, Amazonia, and Baltica. Laurentia lay along the equator.
Recent evidence suggests that South America and Africa never quite joined to Rodinia, though they were located very close to it. Newer reconstructions place Laurentia closer to its present-day orientation, with East Antarctica and Australia to the west, South China to the northwest, Baltica to the east, and Amazonia and Rio de la Plata to the south.
The breakup of Rodinia began by 780 Ma, when numerous mafic dike swarms were emplaced in western Laurentia. Early stages of rifting produced the Belt Supergroup, which is over thick. By 750 Ma the breakup was mostly complete, and Gondwana (composed of most of today's southern continents) had rotated away from Laurentia, which was left isolated near the equator. The breakup of Rodinia may have triggered an episode of severe ice ages (the Snowball Earth hypothesis.)
Pannotia and after
There is some evidence that the fragments of Rodinia gathered into another short-lived supercontinent, Pannotia, at the very end of the Proterozoic. This continent broke up again almost at once, and Laurentia rifted away from South America at around 565 Ma to once again become an isolated continent near the equator, separated from Gondwana by the western Iapetus Ocean. Sometime in the early Cambrian, around 530 Ma, Argentina rifted away from Laurentia and accreted onto Gondwana.
The breakup of Pannotia produced six major continents: Laurentia, Baltica, Kazakhstania, Siberia, China, and Gondwana. Laurentia remained an independent continent until the middle Silurian. During the early to middle Ordovician, several volcanic arcs collided with Laurentia along what is now the Atlantic coast of North America. This caused an episode of mountain-building called the Taconic orogeny. As the mountains raised by the Taconic orogeny were subsequently eroded, they produced the immense Queenston Delta, recorded in the rocks of the Queenston Formation. There was also violent volcanic activity, including the eruption that produced the Millburg/Big Bentonite ash bed. About of ash erupted in this event. However, this does not seem to have triggered any mass extinction.
Throughout the early Paleozoic, Laurentia was characterized by a tectonically stable interior flooded by the seas, with marginal orogenic belts. An important feature was the Transcontinental Arch, which ran southwest from the lowlands of the Canadian Shield. The shield and the arch were the only portions of the continent that were above water through much of the early Paleozoic. There were two major marine transgressions (episodes of continental flooding) during the early Paleozoic, the Sauk and the Tippecanoe. During this time, the Western Cordillera was a passive margin. Sedimentary rocks that were deposited on top of the basement complex were formed in a setting of quiet marine and river waters. The craton was covered by shallow, warm, tropical epicontinental or epicratonic sea (meaning literally "on the craton") that had maximum depths of only about 60 m (200 ft) at the shelf edge.
The position of the equator during the Late Ordovician epoch ( Ma) on Laurentia has been determined via extensive shell bed records. Flooding of the continent that occurred during the Ordovician provided the shallow warm waters for the success of sea life and therefore a spike in the carbonate shells of shellfish. Today the beds are composed of fossilized shells or massive-bedded Thalassinoides facies and loose shells or nonamalgamated brachiopod shell beds. These beds imply the presence of an equatorial climate belt that was hurricane free which lay inside 10° of the equator. This ecological conclusion matches the previous paleomagnetic findings which confirms this equatorial location.
Laurussia
At the end of the Cambrian, about 490 Mya, Avalonia rifted away from Gondwana. By the end of the Ordovician, Avalonia had merged with Baltica, and the two fused to Laurentia at the end of the Silurian (about 420 Ma) in the Caledonian orogeny. This produced the continent of Laurussia.
During this time, several small continental fragments merged with other margins of the craton. These included the North Slope of Alaska, which merged during the Early Devonian. Several small crust fragments accreted from the late Devonian through the Mesozoic to form the Western Cordillera.
The Western Cordillera became a convergent plate margin during the Ordovician, and the Transcontinental Arch became submerged, only to reappear in the Devonian. The Devonian also saw the deposition of the Chattanooga Shale and the Antler Orogeny in the Western Cordillera.
Formation of Pangaea
During the Carboniferous and Permian, Laurussia fused with Gondwana to form Pangaea. The resulting Alleghanian orogeny created the Central Pangean Mountains. The mountains were located close to the equator and produced a year-round zone of heavy precipitation that promoted the deposition of extensive coal beds, including the Appalachian coal beds in the U.S. Meanwhile, Gondwana had drifted onto the South Pole, and cycles of extensive glaciation produced a characteristic pattern of alternating marine and coal swamp beds called cyclothems.
During the Pennsylvanian, the Ancestral Rocky Mountains were raised in the southwestern part of Laurentia. This has been attributed either to either the collision with Gondwana or subduction under the continental margin from the southwest. Two additional marine transgressions took place during the late Paleozoic: the Kaskaskia and Absaroka.
The great continental mass of Pangaea strongly affected climate patterns. The Permian was relatively arid, and evaporites were deposited in the Permian Basin. Sedimentary beds deposited in the southwest in the early Triassic were fluvial in character, but gave way to eolian beds in the late Triassic. Pangaea reached its height about 250 Ma, at the start of the Triassic.
Breakup of Pangaea
The breakup of Pangaea began in the Triassic, with rifting along what is now the east coast of the U.S. that produced red beds, arkosic sandstone, and lake shale deposits. The central Atlantic ocean basin began opening at about 180 Ma. Florida, which had been a part of Gondwana before the assembly of Pangaea, was left with Laurentia during the opening of the central Atlantic. This former Gondwana fragment includes the Carolina Slate belt and parts of Alabama.
The Gulf of Mexico opened during the Late Triassic and Jurassic. This was accompanied by deposition of evaporite beds that later gave rise to salt domes that are important petroleum reservoirs today. Europe rifted away from North America between 140 and 120 Ma, and Laurentia once again became the core of an independent continent with the opening of the North Atlantic in the Paleogene.
Four orogenies occurred in the Mesozoic in the Western Cordillera: the Sonoma, Nevadan, Sevier, and Laramide. The Nevadan orogeny emplaced the extensive batholiths of the Sierra Nevada. The regression of the Sundance Sea in the late Jurassic was accompanied by deposition of the Morrison Formation, notable for its vertebrate fossils.
During Cretaceous times, the Western Interior Seaway ran from the Gulf of Mexico to the Arctic Ocean, dividing North America into eastern and western land masses. From time to time, land masses or mountain chains rose up on the distant edges of the craton and then eroded down, shedding their sand across the landscape. Chalk beds of the Niobrara Formation were deposited at this time, and accretion of crustal fragments continued along the Western Cordillera.
In the Cenozoic
Northeast Mexico was added to the North American craton relatively recently in geological time. This block was formed from the Mesozoic to nearly the present day, with only small fragments of earlier basement rock. It moved as a coherent unit after the breakup of Pangaea. The Atlantic and Gulf Coasts experienced eight transgressions in the Cenozoic. The Laramide orogeny continued to raise the present Rocky Mountains into the Paleocene. The Western Cordillera continued to suffer tectonic deformation, including the formation of the Basin and Range Province in the middle Cenozoic and the uplift of the Colorado Plateau. The Colorado Plateau was uplifted with remarkably little deformation. The flood basalts of the Columbia Plateau also erupted during the Cenozoic.
The southwestern portion of Laurentia consists of Precambrian basement rocks deformed by continental collisions. This area has been subjected to considerable rifting as the Basin and Range Province has been stretched up to 100% of its original width. The area experienced numerous large volcanic eruptions. Baja California rifted away from North America during the Miocene. This block of crust consists of Proterozoic to early Paleozoic shelf and Mesozoic arc volcano formations. The Holocene being an interglacial, a warm spell between episodes of extensive glaciation.
Paleoenvironmental change
Several climate events occurred in Laurentia during the Phanerozoic eon. During the late Cambrian through the Ordovician, sea level fluctuated with ice cap melt. Nine macro scale fluctuations of "global hyper warming", or high intensity greenhouse gas conditions, occurred. Due to sea level fluctuation, these intervals led to mudstone deposits on Laurentia that act as a record of events. The late Ordovician brought a cooling period, although the extent of this cooling is still debated. More than 100 million years later, in the Permian, an overall warming trend occurred. As indicated by fossilized invertebrates, the western margin of Laurentia was affected by a lasting southward bound cool current. This current contrasted with waters warming in the Texas region. This opposition suggests that, during Permian global warm period, northern and northwestern Pangea (western Laurentia) remained relatively cool.
Geological history
Around 4.03 to 3.58 Ga, the oldest intact rock formation on the planet, the Acasta Gneiss, was formed in what is now Northwest Territories (older individual mineral grains are known, but not whole rocks).
Around 2.565 Ga, Arctica formed as an independent continent.
Around 2.72 to 2.45 Ga, Arctica was part of the supercontinent Kenorland.
Around 2.1 to 1.84 Ga, when Kenorland broke apart, the Arctican craton was part of the landmass Nena along with Baltica and Eastern Antarctica.
Around 1.82 Ga, Laurentia was part of the supercontinent Columbia.
Around 1.35–1.3 Ga, Laurentia was an independent continent.
Around 1.3 Ga, Laurentia was part of the landmass Protorodinia.
Around 1.07 Ga, Laurentia was part of the supercontinent Rodinia.
Around 750 Ma, Laurentia was part of the landmass Protolaurasia. Laurentia nearly rifted apart.
In the Ediacaran (635 to 541 ±0.3 Ma), Laurentia was part of the supercontinent Pannotia.
In the Cambrian (541 ±0.3 to 485.4 ±1.7 Ma), Laurentia was an independent continent.
In the Ordovician (485.4 ± 1.7 to 443.8 ±1.5 Ma), Laurentia was shrinking and Baltica was expanding.
In the Devonian (419.2 ± 2.8 to 358.9 ±2.5 Ma), Laurentia collided against Baltica, forming the landmass Euramerica.
In the Permian (298.9 ± 0.8 to 252.17 ±0.4 Ma), all major continents collided against each other, forming the supercontinent Pangaea.
In the Jurassic (201.3 ± 0.6 to 145 ±4 Ma), Pangaea rifted into two landmasses: Laurasia and Gondwana. Laurentia was part of the landmass Laurasia.
In the Cretaceous (145 ± 4 to 66 Ma), Laurentia was an independent continent called North America.
In the Neogene (23.03 ± 0.05 Ma until today or ending 2.588 Ma), Laurentia, in the form of North America, collided with South America, forming the landmass America.
| Physical sciences | Paleogeography | Earth science |
19630739 | https://en.wikipedia.org/wiki/Continent | Continent | A continent is any of several large geographical regions. Continents are generally identified by convention rather than any strict criteria. A continent could be a single landmass or a part of a very large landmass, as in the case of Asia or Europe. Due to this, the number of continents varies; up to seven or as few as four geographical regions are commonly regarded as continents. Most English-speaking countries recognize seven regions as continents. In order from largest to smallest in area, these seven regions are Asia, Africa, North America, South America, Antarctica, Europe, and Australia. Different variations with fewer continents merge some of these regions; examples of this are merging Asia and Europe into Eurasia, North America and South America into America, and Africa, Asia, and Europe into Afro-Eurasia.
Oceanic islands are occasionally grouped with a nearby continent to divide all the world's land into geographical regions. Under this scheme, most of the island countries and territories in the Pacific Ocean are grouped together with the continent of Australia to form the geographical region of Oceania.
In geology, a continent is defined as "one of Earth's major landmasses, including both dry land and continental shelves". The geological continents correspond to seven large areas of continental crust that are found on the tectonic plates, but exclude small continental fragments such as Madagascar that are generally referred to as microcontinents. Continental crust is only known to exist on Earth.
The idea of continental drift gained recognition in the 20th century. It postulates that the current continents formed from the breaking up of a supercontinent (Pangaea) that formed hundreds of millions of years ago.
Etymology
From the 16th century the English noun continent was derived from the term continent land, meaning continuous or connected land and translated from the Latin . The noun was used to mean "a connected or continuous tract of land" or mainland. It was not applied only to very large areas of land—in the 17th century, references were made to the continents (or mainlands) of the Isle of Man, Ireland and Wales and in 1745 to Sumatra. The word continent was used in translating Greek and Latin writings about the three "parts" of the world, although in the original languages no word of exactly the same meaning as continent was used.
While continent was used on the one hand for relatively small areas of continuous land, on the other hand geographers again raised Herodotus's query about why a single large landmass should be divided into separate continents. In the mid-17th century, Peter Heylin wrote in his Cosmographie that "A Continent is a great quantity of Land, not separated by any Sea from the rest of the World, as the whole Continent of Europe, Asia, Africa." In 1727, Ephraim Chambers wrote in his Cyclopædia, "The world is ordinarily divided into two grand continents: the Old and the New." And in his 1752 atlas, Emanuel Bowen defined a continent as "a large space of dry land comprehending many countries all joined together, without any separation by water. Thus Europe, Asia, and Africa is one great continent, as America is another." However, the old idea of Europe, Asia and Africa as "parts" of the world ultimately persisted with these being regarded as separate continents.
Definitions and application
By convention, continents "are understood to be large, continuous, discrete masses of land, ideally separated by expanses of water". By this definition, all continents have to be an island of some metric. In modern schemes with five or more recognized continents, at least one pair of continents is joined by land in some fashion. The criterion "large" leads to arbitrary classification: Greenland, with a surface area of , is only considered the world's largest island, while Australia, at , is deemed the smallest continent.
Earth's major landmasses all have coasts on a single, continuous World Ocean, which is divided into several principal oceanic components by the continents and various geographic criteria.
The geological definition of a continent has four criteria: high elevation relative to the ocean floor; a wide range of igneous, metamorphic and sedimentary rocks rich in silica; a crust thicker than the surrounding oceanic crust; and well-defined limits around a large enough area.
Extent
The most restricted meaning of continent is that of a continuous area of land or mainland, with the coastline and any land boundaries forming the edge of the continent. In this sense, the term continental Europe (sometimes referred to in Britain as "the Continent") is used to refer to mainland Europe, excluding islands such as Great Britain, Iceland, Ireland, and Malta, while the term continent of Australia may refer to the mainland of Australia, excluding New Guinea, Tasmania, and other nearby islands. Similarly, the continental United States refers to "the 49 States (including Alaska but excluding Hawaii) located on the continent of North America, and the District of Columbia."
From the perspective of geology or physical geography, continent may be extended beyond the confines of continuous dry land to include the shallow, submerged adjacent area (the continental shelf) and the islands on the shelf (continental islands), as they are structurally part of the continent.
From this perspective, the edge of the continental shelf is the true edge of the continent, as shorelines vary with changes in sea level. In this sense the islands of Great Britain and Ireland are part of Europe, while Australia and the island of New Guinea together form a continent. Taken to its limit, this view could support the view that there are only three continents: Antarctica, Australia-New Guinea, and a single mega-continent which joins Afro-Eurasia and America via the contiguous continental shelf in and around the Bering Sea. The vast size of the latter compared to the first two might even lead some to say it is the only continent, the others being more comparable to Greenland or New Zealand.
As a cultural construct, the concept of a continent may go beyond the continental shelf to include oceanic islands and continental fragments. In this way, Iceland is considered a part of Europe, and Madagascar a part of Africa. Extrapolating the concept to its extreme, some geographers group the Australian continental landmass with other islands in the Pacific Ocean into Oceania, which is usually considered a region rather than a continent. This divides the entire land surface of Earth into continents, regions, or quasi-continents.
Separation
The criterion that each continent is a discrete landmass is commonly relaxed due to historical conventions and practical use. Of the seven most globally recognized continents, only Antarctica and Australia are completely separated from other continents by the ocean. Several continents are defined not as absolutely distinct bodies but as "more or less discrete masses of land". Africa and Asia are joined by the Isthmus of Suez, and North America and South America by the Isthmus of Panama. In both cases, there is no complete separation of these landmasses by water (disregarding the Suez Canal and the Panama Canal, which are both narrow and shallow, as well as human-made). Both of these isthmuses are very narrow compared to the bulk of the landmasses they unite.
North America and South America are treated as separate continents in the seven-continent model. However, they may also be viewed as a single continent known as America. This viewpoint was common in the United States until World War II, and remains prevalent in some Asian six-continent models. The single American continent model remains a common view in European countries like France, Greece, Hungary, Italy, Malta, Portugal, Spain, Latin American countries and some Asian countries.
The criterion of a discrete landmass is completely disregarded if the continuous landmass of Eurasia is classified as two separate continents (Asia and Europe). Physiographically, Europe and the Indian subcontinent are large peninsulas of the Eurasian landmass. However, Europe is considered a continent with its comparatively large land area of , while the Indian subcontinent, with less than half that area, is considered a subcontinent. The alternative view—in geology and geography—that Eurasia is a single continent results in a six-continent view of the world. Some view the separation of Eurasia into Asia and Europe as a residue of Eurocentrism: "In physical, cultural and historical diversity, China and India are comparable to the entire European landmass, not to a single European country. [...]." However, for historical and cultural reasons, the view of Europe as a separate continent continues in almost all categorizations.
If continents are defined strictly as discrete landmasses, embracing all the contiguous land of a body, then Africa, Asia, and Europe form a single continent which may be referred to as Afro-Eurasia. Combined with the consolidation of the Americas, this would produce a four-continent model consisting of Afro-Eurasia, America, Antarctica, and Australia.
When sea levels were lower during the Pleistocene ice ages, greater areas of the continental shelf were exposed as dry land, forming land bridges between Tasmania and the Australian mainland. At those times, Australia and New Guinea were a single, continuous continent known as Sahul. Likewise, Afro-Eurasia and the Americas were joined by the Bering Land Bridge. Other islands, such as Great Britain, were joined to the mainlands of their continents. At that time, there were just three discrete landmasses in the world: Africa-Eurasia-America, Antarctica, and Australia-New Guinea (Sahul).
Number
There are several ways of distinguishing the continents:
The seven-continent model is taught in most English-speaking countries, including Australia, Canada, the United Kingdom, and the United States, and also in Bangladesh, China, India, Indonesia, Pakistan, the Philippines, Sri Lanka, Suriname, parts of Europe and Africa.
The six-continent combined-Eurasia model is mostly used in Russia and some parts of Eastern Europe.
The six-continent combined-America model is taught in Greece and many Romance-speaking countriesincluding Latin America.
The Olympic flag's five rings represent the five inhabited continents of the combined-America model but excludes the uninhabited Antarctica.
In the English-speaking countries, geographers often use the term Oceania to denote a geographical region which includes most of the island countries and territories in the Pacific Ocean, as well as the continent of Australia.
Eighth continent
Zealandia (a submerged continent) has been called the eighth continent.
Area and population
The following table provides areas given by the Encyclopædia Britannica for each continent in accordance with the seven-continent model, including Australasia along with Melanesia, Micronesia, and Polynesia as parts of Oceania. It also provides populations of continents according to estimates by the United Nations Statistics Division based on the United Nations geoscheme, which includes all of Egypt (including the Isthmus of Suez and the Sinai Peninsula) as a part of Africa, all of Armenia, Azerbaijan, Cyprus, Georgia, Indonesia, Kazakhstan, and Turkey (including East Thrace) as parts of Asia, all of Russia (including Siberia) as a part of Europe, all of Panama and the United States (including Hawaii) as parts of North America, and all of Chile (including Easter Island) as a part of South America.
Other divisions
Supercontinents
Apart from the current continents, the scope and meaning of the term continent includes past geological ones. Supercontinents, largely in evidence earlier in the geological record, are landmasses that comprise most of the world's cratons or continental cores. These have included Vaalbara, Kenorland, Columbia, Rodinia, Pannotia, and Pangaea. Over time, these supercontinents broke apart into large landmasses which formed the present continents.
Subcontinents
Certain parts of continents are recognized as subcontinents, especially the large peninsulas separated from the main continental landmass by geographical features. The most widely recognized example is the Indian subcontinent. The Arabian Peninsula, Southern Africa, the Southern Cone of South America, and Alaska in North America might be considered further examples.
In many of these cases, the "subcontinents" concerned are on different tectonic plates from the rest of the continent, providing a geological justification for the terminology. Greenland, generally considered the world's largest island on the northeastern periphery of the North American Plate, is sometimes referred to as a subcontinent. This is a significant departure from the more conventional view of a subcontinent as comprising a very large peninsula on the fringe of a continent.
Where the Americas are viewed as a single continent (America), it is divided into two subcontinents (North America and South America) or three (Central America being the third). When Eurasia is regarded as a single continent, Asia and Europe are treated as subcontinents.
Submerged continents
Some areas of continental crust are largely covered by the ocean and may be considered submerged continents. Notable examples are Zealandia, emerging from the ocean primarily in New Zealand and New Caledonia, and the almost completely submerged Kerguelen Plateau in the southern Indian Ocean.
Microcontinents
Some islands lie on sections of continental crust that have rifted and drifted apart from a main continental landmass. While not considered continents because of their relatively small size, they may be considered microcontinents. Madagascar, the largest example, is usually considered an island of Africa, but its divergent evolution has caused it to be referred to as "the eighth continent" from a biological perspective.
Geological continents
Geologists use four key attributes to define a continent:
Elevation – The landmass, whether dry or submerged beneath the ocean, should be elevated above the surrounding ocean crust.
Geology – The landmass should contain different types of rock: igneous, metamorphic, and sedimentary.
Crustal structure – The landmass should consist of the continental crust, which is thicker and has a lower seismic velocity than the oceanic crust.
Limits and area – The landmass should have clearly defined boundaries and an area of more than one million square kilometres.
With the addition of Zealandia in 2017, Earth currently has seven recognized geological continents:
Africa
Antarctica
Australia
Eurasia
North America
South America
Zealandia
Due to a seeming lack of Precambrian cratonic rocks, Zealandia's status as a geological continent has been disputed by some geologists. However, a study conducted in 2021 found that part of the submerged continent is indeed Precambrian, twice as old as geologists had previously thought, which is further evidence that supports the idea of Zealandia being a geological continent.
All seven geological continents are spatially isolated by geologic features.
History of the concept
Early concepts of the Old World continents
The term "continent" translates the Greek word , meaning "landmass, terra firma", the proper name of Epirus and later especially used for Asia (i.e. Asia Minor).
The first distinction between continents was made by ancient Greek mariners who gave the names Europe and Asia to the lands on either side of the waterways of the Aegean Sea, the Dardanelles strait, the Sea of Marmara, the Bosporus strait and the Black Sea. The names were first applied just to lands near the coast and only later extended to include the hinterlands. But the division was only carried through to the end of navigable waterways and "... beyond that point the Hellenic geographers never succeeded in laying their finger on any inland feature in the physical landscape that could offer any convincing line for partitioning an indivisible Eurasia ..."
Ancient Greek thinkers subsequently debated whether Africa (then called Libya) should be considered part of Asia or a third part of the world. Division into three parts eventually came to predominate. From the Greek viewpoint, the Aegean Sea was the center of the world; Asia lay to the east, Europe to the north and west, and Africa to the south. The boundaries between the continents were not fixed. Early on, the Europe–Asia boundary was taken to run from the Black Sea along the Rioni River (known then as the Phasis) in Georgia. Later it was viewed as running from the Black Sea through Kerch Strait, the Sea of Azov and along the Don River (known then as the Tanais) in Russia. The boundary between Asia and Africa was generally taken to be the Nile River. Herodotus in the 5th century BCE objected to the whole of Egypt being split between Asia and Africa ("Libya") and took the boundary to lie along the western border of Egypt, regarding Egypt as part of Asia. He also questioned the division into three of what is really a single landmass, a debate that continues nearly two and a half millennia later. Herodotus believed Europe to be larger (at least in width) than the other two continents:
I wonder, then, at those who have mapped out and divided the world into Libya, Asia, and Europe; for the difference between them is great, seeing that in length Europe stretches along both the others together, and it appears to me to be wider beyond all comparison.
Eratosthenes, in the 3rd century BCE, noted that some geographers divided the continents by rivers (the Nile and the Don), thus considering them "islands". Others divided the continents by isthmuses, calling the continents "peninsulas". These latter geographers set the border between Europe and Asia at the isthmus between the Black Sea and the Caspian Sea, and the border between Asia and Africa at the isthmus between the Red Sea and the mouth of Lake Bardawil on the Mediterranean Sea.
The Roman author Pliny the Elder, writing in the 1st century CE, stated that "The whole globe is divided into three parts, Europe, Asia, and Africa", adding:
I shall first then speak of Europe, the foster-mother of that people which has conquered all other nations, and itself by far the most beauteous portion of the earth. Indeed, many persons have, not without reason, considered it, not as a third part only of the earth, but as equal to all the rest, looking upon the whole of our globe as divided into two parts only, by a line drawn from the river Tanais to the Straits of Gades.
Following the fall of the Western Roman Empire, the culture that developed in its place, linked to Latin and the Catholic church, began to associate itself with the concept of Europe. Through the Roman period and the Middle Ages, a few writers took the Isthmus of Suez as the boundary between Asia and Africa, but most writers continued to consider it the Nile or the western border of Egypt (Gibbon). In the Middle Ages, the world was usually portrayed on T and O maps, with the T representing the waters dividing the three continents. By the middle of the 18th century, "the fashion of dividing Asia and Africa at the Nile, or at the Great Catabathmus [the boundary between Egypt and Libya] farther west, had even then scarcely passed away".
European arrival in the Americas
Christopher Columbus sailed across the Atlantic Ocean to the Caribbean in 1492, sparking a period of European exploration of the Americas. But despite four voyages to the Americas, Columbus never believed he had reached a new continent—he always thought it was part of Asia.
In 1501, Amerigo Vespucci and Gonçalo Coelho attempted to sail around what they considered the southern end of the Asian mainland into the Indian Ocean, passing through Fernando de Noronha. After reaching the coast of Brazil, they sailed along the coast of South America much farther south than Asia was known to extend, confirming that this was a land of continental proportions. On return to Europe, an account of the voyage, called Mundus Novus ("New World"), was published under Vespucci's name in 1502 or 1503, although it seems that it had additions or alterations by another writer. Regardless of who penned the words, Mundus Novus credited Vespucci with saying, "I have discovered a continent in those southern regions that is inhabited by more numerous people and animals than our Europe, or Asia or Africa", the first known explicit identification of part of the Americas as a continent like the other three.
Within a few years, the name "New World" began appearing as a name for South America on world maps, such as the Oliveriana (Pesaro) map of around 1504–1505. Maps of this time, though, still showed North America connected to Asia and showed South America as a separate land.
In 1507 Martin Waldseemüller published a world map, Universalis Cosmographia, which was the first to show North and South America as separate from Asia and surrounded by water. A small inset map above the main map explicitly showed for the first time the Americas being east of Asia and separated from Asia by an ocean, as opposed to just placing the Americas on the left end of the map and Asia on the right end. In the accompanying book Cosmographiae Introductio, Waldseemüller noted that the earth is divided into four parts, Europe, Asia, Africa, and the fourth part, which he named "America" after Amerigo Vespucci's first name. On the map, the word "America" was placed on part of South America.
Beyond four continents
The Sanskrit text Rig Veda often dated 1500 BCE has the earliest mention of seven continents in the Earth, the text claims that the Earth has seven continents and Lord Vishnu Measured the entire universe from his first foot from the land of Earth which has 7 continents.
In regard to the above-quoted verses, it is commonly accepted that there are Seven Continents or 'regions of the earth'. A. Glucklich adds that 'In the Matsya Purana, for instance, there is a seven-part map of the world ... [it has] one centre, where an immense mountain – Mount Meru (or Maha Meru, Great Meru) – stands ... The continents encircle the mountain in seven concentric circles ... It seems clear that the Himalayas were the approximate location of Mt. Meru and the text is clear that the earth has seven continents.
From the late 18th century, some geographers started to regard North America and South America as two parts of the world, making five parts in total. Overall though, the fourfold division prevailed well into the 19th century.
Europeans discovered Australia in 1606, but for some time it was taken as part of Asia. By the late 18th century, some geographers considered it a continent in its own right, making it the sixth (or fifth for those still taking America as a single continent). In 1813, Samuel Butler wrote of Australia as "New Holland, an immense island, which some geographers dignify with the appellation of another continent" and the Oxford English Dictionary was just as equivocal some decades later. It was in the 1950s that the concept of Oceania as a "great division" of the world was replaced by the concept of Australia as a continent.
Antarctica was sighted in 1820 during the First Russian Antarctic Expedition and described as a continent by Charles Wilkes on the United States Exploring Expedition in 1838, the last continent identified, although a great "Antarctic" (antipodean) landmass had been anticipated for millennia. An 1849 atlas labelled Antarctica as a continent but few atlases did so until after World War II.
Over time, the western concept of dividing the world into continents spread globally, replacing conceptions in other areas of the world. The idea of continents continued to become imbued with cultural and political meaning. In the 19th century during the Meiji period, Japanese leaders began to self-identify with the concept of being Asian, and renew relations with other "Asian" countries while conceiving of the idea of Asian solidarity against western countries. This conception of an Asian identity, as well as the idea of Asian solidarity, was later taken up by others in the region, such as Republican China and Vietnam.
From the mid-19th century, atlases published in the United States more commonly treated North and South America as separate continents, while atlases published in Europe usually considered them one continent. However, it was still not uncommon for American atlases to treat them as one continent up until World War II. From the 1950s, most U.S. geographers divided the Americas into two continents. With the addition of Antarctica, this made the seven-continent model. However, this division of the Americas never appealed to Latin Americans, who saw their region spanning an as a single landmass, and there the conception of six continents remains dominant, as it does in scattered other countries.
Some geographers regard Europe and Asia together as a single continent, dubbed Eurasia. In this model, the world is divided into six continents, with North America and South America considered separate continents.
Geology
Geologists use the term continent in a different manner from geographers. In geology, a continent is defined by continental crust, which is a platform of metamorphic and igneous rocks, largely of granitic composition. Continental crust is less dense and much thicker than oceanic crust, which causes it to "float" higher than oceanic crust on the dense underlying mantle. This explains why the continents form high platforms surrounded by deep ocean basins.
Some geologists restrict the term continent to portions of the crust built around stable regions called cratons. Cratons have largely been unaffected by mountain-building events (orogenies) since the Precambrian. A craton typically consists of a continental shield surrounded by a continental platform. The shield is a region where ancient crystalline basement rock (typically 1.5 to 3.8 billion years old) is widely exposed at the surface. The platform surrounding the shield is also composed of ancient basement rock, but with a cover of younger sedimentary rock. The continents are accretionary crustal "rafts" that, unlike the denser basaltic crust of the ocean basins, are not subjected to destruction through the plate tectonic process of subduction. This accounts for the great age of the rocks comprising the continental cratons.
The margins of geologic continents are either active or passive. An active margin is characterised by mountain building, either through a continent-on continent collision or a subduction zone. Continents grow by accreting lighter volcanic island chains and microcontinents along these active margins, forming orogens. At a passive margin, the continental crust is stretched thin by extension to form a continental shelf, which tapers off with a gradual slope covered in sediment, connecting it directly to the oceanic crust beyond. Most passive margins eventually transition into active margins: where the oceanic plate becomes too heavy due to cooling, it disconnects from the continental crust, and starts subducting below it, forming a new subduction zone.
There are many microcontinents, or continental fragments, that are built of continental crust but do not contain a craton. Some of these are fragments of Gondwana or other ancient cratonic continents: Zealandia, which includes New Zealand and New Caledonia; Madagascar; the northern Mascarene Plateau, which includes the Seychelles. Other islands, such as several in the Caribbean Sea, are composed largely of granitic rock as well, but all continents contain both granitic and basaltic crust, and there is no clear boundary as to which islands would be considered microcontinents under such a definition. The Kerguelen Plateau, for example, is largely volcanic, but is associated with the breakup of Gondwanaland and is considered a microcontinent, whereas volcanic Iceland and Hawaii are not. The British Isles, Sri Lanka, Borneo, and Newfoundland were on the margins of the Laurasian continent—only separated from the main continental landmass by inland seas flooding its margins.
The movement of plates has caused the continual formation and breakup of continents, and occasionally supercontinents, in a process called the Wilson Cycle. The supercontinent Columbia or Nuna formed during a period of 2.0–1.8 billion years ago and broke up about 1.5–1.3 billion years ago. The supercontinent Rodinia is thought to have formed about 1 billion years ago and to have embodied most or all of Earth's continents, and broken up into eight continents around 600 million years ago. The eight continents later reassembled into another supercontinent called Pangaea; Pangaea broke up into Laurasia (which became North America and Eurasia) and Gondwana (which became the remaining continents).
Criticism
Some academics, such as the historical geographer Martin W. Lewis, argue that the systems we understand today are more rooted in social, political, and cultural history than in geological fact, a view particularly outlined in his book The Myth of Continents: A Critique of Metageography.
| Physical sciences | Continents and regions | null |
23649300 | https://en.wikipedia.org/wiki/Opuntia | Opuntia | Opuntia, commonly called the prickly pear cactus, is a genus of flowering plants in the cactus family Cactaceae, many known for their flavorful fruit and showy flowers. Cacti are well-adapted to aridity; however, they are still vulnerable to alterations in precipitation and temperature driven by climate change. Prickly pear alone is more commonly used to refer exclusively to the fruit, but may also be used for the plant itself; in addition, other names given to the plant and its specific parts include tuna (fruit), sabra, sabbar, nopal (pads, plural nopales) from the Nahuatl word , nostle (fruit) from the Nahuatl word , and paddle cactus. The genus is named for the Ancient Greek city of Opus, where, according to Theophrastus, an edible plant grew and could be propagated by rooting its leaves. The most common culinary species is the "Barbary fig" (Opuntia ficus-indica).
Description
O. ficus-indica is a large, trunk-forming, segmented cactus that may grow to with a crown of over in diameter and a trunk diameter of . Cladodes (large pads) are green to blue-green, bearing few spines up to or may be spineless. Prickly pears typically grow with flat, rounded cladodes (also called platyclades) containing large, smooth, fixed spines and small, hairlike prickles called glochids that readily adhere to skin or hair, then detach from the plant. The flowers are typically large, axillary, solitary, bisexual, and epiperigynous, with a perianth consisting of distinct, spirally arranged tepals and a hypanthium. The stamens are numerous and in spiral or whorled clusters, and the gynoecium has numerous inferior ovaries per carpel. Placentation is parietal, and the fruit is a berry with arillate seeds. Prickly pear species can vary greatly in habit; most are shrubs, but some, such as O. galapageia of the Galápagos, are trees.
Growth
Chemistry
Opuntia contains a range of phytochemicals in variable quantities, such as polyphenols, dietary minerals and betalains. Identified compounds under basic research include gallic acid, vanillic acid and catechins, as examples. The Sicilian prickly pear contains betalain, betanin, and indicaxanthin, with highest levels in their fruits.
Taxonomy
[[File:Opuntia Cactus in Behbahan.jpg|alt=O. lindheimeri blooming, Behbahan|thumb|O. lindheimeri, Behbahan]]
When Carl Linnaeus published Species Plantarum in 1753 – the starting point for modern botanical nomenclature – he placed all the species of cactus known to him in one genus, Cactus. In 1754, the Scottish botanist Philip Miller divided them into several genera, including Opuntia. He distinguished the genus largely on the form of its flowers and fruits.
Considerable variation of taxonomy occurs within Opuntia species, resulting in names being created for variants or subtypes within a species, and use of DNA sequencing to define and isolate various species.
SpeciesSee List of Opuntia speciesOpuntia hybridizes readily between species. This can make classification difficult, yielding a reticulate phylogeny where different species come together in hybridization. Opuntia also has a tendency for polyploidy. The ancestral diploid state was 2n=22, but many species are hexaploid (6n = 66) or octaploid (8n = 88).
Formerly in Opuntia
Austrocylindropuntia Brasiliopuntia Corynopuntia Cylindropuntia Disocactus phyllanthoides (as O. speciosa)
Micropuntia MiqueliopuntiaChollas
Chollas, now recognized to belong to the distinct genus Cylindropuntia, are distinguished by having cylindrical, rather than flattened, stem segments with large barbed spines. The stem joints of several species, notably the jumping cholla (C. fulgida), are very brittle on young stems, readily breaking off when the barbed spines stick to clothing or animal fur as a method of vegetative reproduction. The barbed spines can remain embedded in the skin, causing discomfort and sometimes injury.
Breeding
One of the ancient homes of the cactus pear, Mexico, ran a breeding program in the 1960s. This effort at the Antonio Narro Agrarian Autonomous University (Universidad Autónoma Agraria Antonio Narro, UAAAN) produced improvements in some traits including cold-hardiness.
Distribution and habitat
Like most true cactus species, prickly pears are native only to the Americas. Through human action, they have since been introduced to many other areas of the world. Prickly pear species are found in abundance in Mexico, especially in the central and western regions, and in the Caribbean islands (West Indies). In the United States, prickly pears are native to many areas of the arid, semi-arid, and drought-prone Western and South Central United States, including the lower elevations of the Rocky Mountains and southern Great Plains, where species such as O. phaeacantha and O. polyacantha have become dominant, and to the desert Southwest, where several types are endemic. Prickly pear cactus is also native to sandy coastal beach scrub environments of the East Coast from Florida to southern Connecticut, where O. humifusa, O. stricta, and O. pusilla, are found from the East Coast south into the Caribbean and the Bahamas. Additionally, the eastern prickly pear is native to the midwestern "sand prairies" near major river systems, such as the Mississippi, Illinois, and Ohio rivers. The plant also occurs naturally in hilly areas of southern Illinois, and sandy or rocky areas of northern Illinois.Opuntia species are the most cold-tolerant of the lowland cacti, extending into western and southern Canada. One subspecies, O. fragilis var. fragilis, has been found growing along the Beatton River in north-eastern British Columbia, southwest of Cecil Lake at 56° 17’ N latitude and 120° 39’ W longitude. Others are seen in the Kleskun Hills Natural Area of north-west Alberta at 55° 15’ 30’’ N latitude and 118° 30’ 36’’ W longitude.
Prickly pears produce a fruit known as tuna, commonly eaten in Mexico and in the Mediterranean region, which is also used to make aguas frescas. The fruit can be red, wine-red, green, or yellow-orange. In the Galápagos Islands, the Galápagos prickly pear, O. galapageia, has previously been treated as a number of different species, but is now only divided into varieties and subvarieties. Most of these are confined to one or a few islands, so they have been described as "an excellent example of adaptive radiation". On the whole, islands with tall, trunked varieties are also the home of giant tortoises, whereas islands lacking tortoises have low or prostrate forms of Opuntia. Prickly pears are a prime source of food for the common giant tortoises in the Galápagos Islands, so they are important in the food web.
Charles Darwin was the first to note that the cacti have thigmotactic anthers. When the anthers are touched, they curl over, depositing their pollen on the pollinator. That movement can be seen by gently poking the anthers of an open Opuntia flower. The same trait has evolved convergently in other genera (e.g. Lophophora).
Prickly pears (mostly O. stricta) were originally imported into Europe during the 16th century. They are now found in the Mediterranean region of Northern Africa, especially in Algeria, Morocco and Tunisia, where they grow all over the countryside, and in parts of Southern Europe, especially Spain, where they can be found in the east, south-east, and south of the country, and also in Malta, where they grow all over the islands, and in southern Italy, especially in Sicily and Sardinia. They can be found in enormous numbers in parts of South Africa, where they were introduced from South America. The prickly pear is considered an invasive species in Australia, Ethiopia, South Africa, and Hawaii, among other locations.
The first introduction of prickly pears into Australia is ascribed to the founding governor of the New South Wales colony, Arthur Phillip, and the earliest European colonists, in 1788. Brought from Brazil to Sydney, they were most likely O. monacantha. That variety did not spread beyond the east coast. However, a number of other types of prickly pear were introduced to Australian gardens in the mid-19th century. The cactus was also used as agricultural fencing and a feedstock for animals in times of drought,Patterson, Ewen K. 1936. The World's First Insect Memorial. "The Review of the River Plate", December pp. 16–17 as well as in an attempt to establish a cochineal dye industry.
The cactus quickly became a widespread invasive weed in the dry interior climate west of the Great Dividing Range, in New South Wales and Queensland, eventually converting of farming land into an impenetrable green jungle of prickly pear, in places high. Scores of farmers were driven off their land by what they called the "green hell", and their abandoned homes were crushed under the cactus growth, which advanced at a rate of per year.
In 1919, the Australian federal government established the Commonwealth Prickly Pear Board to coordinate efforts with state governments to eradicate the weed. Early attempts, comprising mechanical removal and poisonous chemicals failed. As a last resort, biological control was attempted. In 1925, the Cactoblastis cactorum moth was introduced from South America, and its larvae rapidly began to control the infestation. Alan Dodd, the son of the noted entomologist Frederick Parkhurst Dodd, was a leading official in combating the prickly pear menace. A memorial hall in Boonarga, Queensland, commemorates the efforts of the moth. The release of cochineal insects, which eat the cactus and simultaneously kill the plant, has also proven an effective measure for combating its spread.
Natural distribution of the plant occurs via consumption and associated seed dispersal by many animals, including antelopes, nonhuman primates, elephants, birds, and humans. When ingested by elephants, the sharp components of the plant cause harm to the mouth, stomach, and intestines.
Ecology O. ficus-indica thrives in regions with mild winters having a prolonged dry spell followed by hot summers with occasional rain and relatively low humidity. A mean annual rainfall of provides good growth rates. O. ficus-indica proliferates in various soils ranging from subacid to subalkaline, with clay content not exceeding 15–20% and the soil well-drained. The shallow root system enables the plant to grow in shallow, loose soils, such as on mountain slopes. Opuntia spreads into large clonal colonies, which contribute to its being considered a noxious weed in some places.Opuntia species are primarily pollinated by bees, including some bee genera (Diadasia and Lithurgus) that contain specialist pollinators (oligoleges) that exclusively visit Opuntia. Only a few Opuntia species, such as O. cochenillifera and O. stenopetala, are pollinated by hummingbirds.
Animals that eat Opuntia include the prickly pear island snail and Cyclura rock iguanas. The fruit are relished by many arid-land animals, chiefly birds, which thus help distribute the seeds. Opuntia pathogens include the sac fungus Colletotrichum coccodes and Sammons' Opuntia virus. The ant Crematogaster opuntiae and the spider Theridion opuntia are named because of their association with prickly pear cactus.
Toxicity
Although the plants are edible, the pointed hairs should not be eaten, and similar species with milky sap are suspect.
Uses
Nutrition
Raw opuntia leaves are 88% water, 10% carbohydrates, and less than 1% both of protein and fat. In a reference serving, raw leaves provide of food energy, 17% of the Daily Value (DV) for vitamin C, and 24% DV for magnesium, with no other micronutrients in significant content.
Regional food uses
The fruit of prickly pears, commonly called cactus fruit, cactus fig, Indian fig (meaning "Native American", not "of India"), nopales or tuna in Spanish, is edible, although it must be peeled carefully to remove the small spines on the outer skin before consumption. If the outer layer is not properly removed, glochids can be ingested, causing discomfort of the throat, lips, and tongue, as the small spines are easily lodged in the skin. Native Americans like the Tequesta would roll the fruit around in a suitable medium (e.g. grit) to "sand" off the glochids. Alternatively, rotating the fruit in the flame of a campfire or torch has been used to remove the glochids. Today, parthenocarpic (seedless) cultivars are also available. The seeds can be used for flour.
In Mexico, prickly pears are often used to make appetizers, soups, salads, entrees, vegetable dishes, breads, desserts, beverages, candy, jelly, and drinks. The young stem segments, usually called pads or nopales, are also edible in most species of Opuntia. They are commonly used in Mexican cuisine in dishes such as huevos con nopales (eggs with nopal), or tacos de nopales. Nopales are also an important ingredient in New Mexican cuisine. In 2009 it was introduced as a cheaper alternative to corn for the production of tortillas and other corn products. They can also be pickled.Opuntia ficus-indica has been introduced to Europe, and flourishes in areas with a suitable climate, such as the south of France and southern Italy: In Sicily, they are referred to as fichi d'India (Italian literal translation of Indian fig) or ficurinia (Sicilian language literal translation of Indian fig). In Sardinia, they are called figumorisca ("Moorish figs"), the same denomination they receive along the Catalan-speaking regions of the Western Mediterranean, figa de moro. They can be found also in the Struma River in Bulgaria, in southern Portugal and Madeira (where they are called tabaibo, figo tuno, or "Indian figs"), in Andalusia, Spain (where they are known as higos chumbos).
In Greece, it grows in such places as the Peloponnese region, Ionian Islands, or Crete, and its figs are known as frangosyka (Frankish, i.e. Western European, figs) or pavlosyka ("Paul's figs"), depending on the region. In Albania, they are called fiq deti translated as "sea figs", and are present in the south-west shore. The figs are also grown in Cyprus, where they are known as papoutsósyka or babutsa ("shoe figs").
The prickly pear also grows widely on the islands of Malta, where it is enjoyed by the Maltese as a typical summer fruit (known as bajtar tax-xewk, literally "spiny figs"), as well as being used to make the popular liqueur known as bajtra. The prickly pear is so commonly found in the Maltese islands, it is often used as a dividing wall between many of Malta's characteristic terraced fields in place of the usual rubble walls.
The prickly pear was introduced to Eritrea during the period of Italian colonisation between 1890 and 1940. It is locally known there as beles and is abundant during the late summer and early autumn (late July through September). The beles from the holy monastery of Debre Bizen is said to be particularly sweet and juicy.
In Morocco, Tunisia, Libya, Saudi Arabia, Jordan, and other parts of North Africa and the Middle East, prickly pears of the yellow and orange varieties are grown by the side of farms, beside railway tracks and other otherwise noncultivable land. It is sold in summer by street vendors, and is considered a refreshing fruit for that season. In Libya, it is a popular summer fruit and called by the locals hindi, which literally means "Indian".Tungi is the local St. Helenian name for cactus pears. The plants (Indian fig opuntia) were originally brought to the island by the colonial ivory traders from East Africa in the 1850s. Tungi cactus now grows wild in the dry coastal regions of the island. Three principal cultivars of tungi grow on the island: the "English" with yellow fruit; the "Madeira" with large red fruit; and the small, firm "spiny red". Tungi also gives its name to a local Spirit distilled at The St Helena distillery at Alarm Forest, the most remote distillery in the world, made entirely from the opuntia cactus.
Cactus pear is being promoted and researched by ICARDA for India, Jordan, and Pakistan especially. It is an underappreciated crop in these countries and has undergone recent expansion in cultivated area. In some particularly promising areas of India and Pakistan it has given a 30% increase in milk yield /hectare (/acre).
Folk medicine
In Mexican folk medicine, its pulp and juice are considered treatments for wounds and inflammation of the digestive and urinary tracts, although there is no high-quality evidence for any clinical benefit of using opuntia for these purposes.
Prior to modern medicine, Native Americans and Mexicans primarily used Opuntia as a coagulant for open wounds, using the pulp of the stem either by splitting the stem or scraping out the pulp.
In one recent study, it was found that Opuntia aided in the prevention or slow down of diabetes, obesity, metabolic syndrome, cardiovascular disease, and some cancers. The results of the group that was taking Opuntia showed a reduction in BMI, body composition, and waist circumference when compared to the placebo group.
Other uses
In dye production Dactylopius coccus is a scale insect from which cochineal dye is derived. D. coccus itself is native to tropical and subtropical South America and Mexico. This insect, a primarily sessile parasite, lives on cacti from the genus Opuntia, feeding on moisture and nutrients in the cactus sap. The insect produces carminic acid, which deters predation by other insects. The carminic acid can be extracted from the insect's body and eggs to make the red dye.
Cochineal is used primarily as a red food colouring and for cosmetics. The cochineal dye was used by the Aztec and Maya peoples of Central and North America, and by the Inca in South America. Produced almost exclusively in Oaxaca, Mexico, by indigenous producers, cochineal became Mexico's second-most valued export after silver. The dyestuff was consumed throughout Europe, and was so highly valued, its price was regularly quoted on the London and Amsterdam Commodity Exchanges.
The biggest producers of cochineal are Peru, the Canary Islands, and Chile. Current health concerns over artificial food additives have renewed the popularity of cochineal dyes, and the increased demand is making cultivation for insect farming an attractive opportunity in other regions, such as in Mexico, where cochineal production had declined again owing to the numerous natural enemies of the scale insect.
Apart from cochineal, the red dye betanin can be extracted from some Opuntia plants themselves. The Navajo have traditionally produced a reddish dye from the fruit of the prickly pear cactus, used in dyeing woolen yarns, and where, after pulverizing the fruit, they are placed in a bath solution of cold water for two or three weeks.
For animal fodder
Cactus is used as a fodder crop for animals in arid and dryland regions. Some farmers prepare it with a fermentation produce, to remove the spines, and increase the digestibility.
As a source of "vegan leather"
The thick skin of nopal cactus can be harvested as an environmentally-friendly leather replacement.
For fuel
Bioethanol can be produced from some Opuntia species.
For bioplastic
Nopal juice can be used to produce bioplastic.
Culture
The prickly pear cactus has been used for centuries both as a food source and a natural fence that keeps in livestock and marks the boundaries of family lands. They are resilient and often grow back following removal.
The 1975–1988 version of the emblem of Malta also featured a prickly pear, along with a traditional dgħajsa, a shovel and pitchfork, and the rising sun.
The prickly pear is the official plant of Texas by legislation from 1995.
The cactus lends its name to a song by British jazz/classical group Portico Quartet. The song "My Rival", on the album Gaucho by the American jazz-pop group Steely Dan begins with the words, "The wind was driving in my face/The smell of prickly pear."
In the fall of 1961, Cuba had its troops plant a barrier of Opuntia cactus along the northeastern section of the fence surrounding the Guantanamo Bay Naval Base to stop Cubans from escaping Cuba to take refuge in the United States. This was dubbed the "Cactus Curtain", an allusion to Europe's Iron Curtain and the Bamboo Curtain in East Asia.
Uruguayan-born footballer Bruno Fornaroli is nicknamed prickly pear due to his sometimes spiky hairstyles.
Mexico
The coat of arms of Mexico depicts a Mexican golden eagle, perched upon an Opuntia cactus, holding a rattlesnake. According to the official history of Mexico, the coat of arms is inspired by an Aztec legend regarding the founding of Tenochtitlan. The Aztecs, then a nomadic tribe, were wandering throughout Mexico in search of a divine sign to indicate the precise spot upon which they were to build their capital. Their god Huitzilopochtli had commanded them to find an eagle devouring a snake, perched atop a cactus that grew on a rock submerged in a lake. After 200 years of wandering, they found the promised sign on a small island in the swampy Lake Texcoco. There they founded their new capital, Tenochtitlan. The cactus (O. ficus-indica; Nahuatl: tenochtli), full of fruits, is the symbol for the island of Tenochtitlan.
Israeli-born Jews
The cactus fig is called tzabar''' in Hebrew (). This cactus is also the origin of the term sabra used to describe any Jew born in Israel. The allusion is to a thorny, spiky skin on the outside, but a soft, sweet interior, suggesting, though the Israeli sabras are rough on the outside, they are sweet and sensitive once one gets to know them. This term is derived from an Arabic word for this cactus صبار ṣubbār, where the related term sabr also translates to "patience" or "tenacity".
Palestinians
The prickly pear is also considered a national symbol of Palestine, having been grown across historic Palestine for years, traditionally being used to mark out land boundaries. The plant is seen by Palestinians as representing qualities of resilience and patience, as represented by the Palestinian proverb saber as-sabbar'' ("the patience of the cactus"). Its use as an emblem of Palestine has been traced to a painting produced by the artist Zulfa al-Sa'di in the 1930s. Its visibility was renewed by historical research carried out in the 1980s and 1990s on the Nakba, which revealed that many destroyed Palestinian villages subsequently saw regrowth of the cacti. The plant's ability to thrive anywhere is also considered to reflect the experiences of the Palestinian diaspora.
| Biology and health sciences | Other culinary fruits | Plants |
3531757 | https://en.wikipedia.org/wiki/King%20quail | King quail | The king quail (Synoicus chinensis), also known as the blue-breasted quail, Asian blue quail, Chinese painted quail, or Chung-Chi, is a species of Old World quail in the family Phasianidae. This species is the smallest "true quail", ranging in the wild from southern China, South and Southeast Asia to Oceania, south to southeastern Australia, with 9 different subspecies. A failed attempt was made to introduce this species to New Zealand by the Otago Acclimatisation Society in the late 1890s. It is quite common in aviculture worldwide, where it is sometimes misleadingly known as the "button quail", which is the name of an only very distantly related family of birds, the buttonquails.
Description
Male king quail occur in many colors, including blue, brown, silver, maroon, dark brown, and almost black. They have orange feet that are hard and able to withstand a continuous life on the ground like many other game birds. Females are similar to the males, but do not occur in shades of blue. They can live up to 13 years in captivity but typically only 3–6 years. In the wild, they may live only 1.5 years. The eggs of king quail are a light, creamy-brown colour and slightly pointed at the "top", roughly ovular in shape. Male quail give a descending whistle and a raspy "snoring" call.
Taxonomy
Six subspecies are recognized:
S. c. chinensis (Linnaeus, 1766): Found from India and Sri Lanka to Malaya, Indochina, southeastern China, and Taiwan
S. c. trinkutensis Richmond, 1902: Nicobar blue-breasted quail, found on the Nicobar Islands
S. c. lineatus (Scopoli, 1786): Found in the Philippines, Borneo, Lesser Sundas, Sulawesi and Sula Islands
S. c. lepidus (Hartlaub, 1879): Found in New Guinea and the Bismarck Archipelago
S. c. victoriae (Mathews, 1912): Found in eastern Australia
S. c. colletti (Mathews, 1912): Found in northern Australia
The species has had a complex taxonomic history, being classified into the genus Coturnix, then Synoicus, then Excalfactoria. Phylogenetic evidence supports it belonging in an expanded Synoicus that, alongside the blue quail (S. adansonii) also includes the Snow Mountain quail (S. monorthonyx) and brown quail (S. ypsilophorus).
The subspecies S. c. victoriae was formerly named C. s. australis (Gould, 1865), but the reclassification of the species into Synoicus caused this designation to be preoccupied by the Australian subspecies of the brown quail (S. y. australis), leading to the epithet being changed to victoriae (Mathews, 1912).
Reproduction
The males fight for the right to mate with the females. The winner then breeds every female. Females can then develop and lay an egg within one to two days of being bred. They either build a nest first or lay eggs on the ground. Females usually only go broody when they have collected an ideal clutch size. Clutch size varies from five to 13 eggs. Before incubation starts, all the eggs composing the clutch will have been laid. In captivity, the ideal number of eggs in a clutch is six to eight. The chicks hatch after about 16 days.
Conservation status
Australia
King quail are not listed as threatened on the Australian Environment Protection and Biodiversity Conservation Act 1999.
State of Victoria, Australia
This species is listed as threatened on the Victorian Flora and Fauna Guarantee Act (1988). Under this Act, an Action Statement for the recovery and future management of this species has not been prepared.
On the 2007 advisory list of threatened vertebrate fauna in Victoria, this species is listed as endangered.
Aviculture
This quail has been very popular to keep and breed for many years; numerous mutations have been developed. They are quite hardy once they have adjusted to their surroundings and keep the bottom of an aviary spotless. A great advantage of these quail is that they live exclusively on the ground, and do not interfere with other birds. The cost of purchasing and maintaining them is very little. They have been known to become hand-tame.
They may be housed in pairs to quartets in a planted aviary, kept singly in bird cages, or in colonies in large flights. Males may compete, as may females. Suspension cages do not work well for this species of quail because of their smaller feet; a much finer size of floor wire should be employed.
Females lay an egg a day if kept on the proper diet. Nesting sites can be as spartan as a quiet corner or a depression in the ground against a wall. Preferably, a clump of long grass, tea tree branches, or pile of loose herbage should be provided. Often, a hen lays eggs on the aviary floor without the use of a nest. This is a sign that the birds are not content with the existing facilities and the provision of a sheltered nest site may result in a nest being built. The cock usually selects the nest site. The nest is a simple scrape in the ground, lined with grasses, and is built by the hen with some assistance from the cock. The eggs measuring 25 x 19mm are variable in colour from the palest of browns to dark olive and peppered with fine black spots. Clutch size varies from four to 13, but occasionally a hen can be found incubating upwards of 20 eggs. It is usually a combined clutch from a number of hens, and due to the difficulties of turning and covering a clutch of that size, hatching rate is often poor. Removing some of the eggs and artificially incubating or fostering them may be beneficial.
The species usually breeds year-round; incubation times are from 18 to 23 days before chicks hatch. The hen cares for the chicks until around 4 weeks of age, when they should be separated from parent birds into a separate aviary.
Hybrids and mutations
Hybrids of king quail and brown quail are known.
Silvers and cinnamon are the most common colour varieties. Pied, albino, and charcoals are becoming more common. Mutations can be combined.
Occasionally, cock-feathered hens appear; this is not a mutation as such, but one of a few conditions that affect normal hormonal balances. It is most often seen when a hen has an ovarian cyst or growth. They usually stop laying eggs, but can live for a number of years happily just looking like a male. In one case, a silver hen was kept for many years by herself, moulted into cock plumage, and laid only extremely pale, green-shelled eggs for a few seasons before passing of old age.
Diet
In the wild, the diet of king quail consists of small bugs, grass seeds, rape seeds, and vegetation that is available at the time. In aviculture, all birds should be fed a variety of seeds and a healthy range of fruit and vegetables. During breeding, hens should be fed calcium-rich food sources such as shell grit to prevent egg binding. Newly hatched chicks should be fed high-protein chick crumb mixed with a little water. Other sources of protein include mealworms, termites and various insects.
| Biology and health sciences | Galliformes | Animals |
3533321 | https://en.wikipedia.org/wiki/Amargasaurus | Amargasaurus | Amargasaurus (; "La Amarga lizard") is a genus of sauropod dinosaur from the Early Cretaceous epoch (129.4–122.46 mya) of what is now Argentina. The only known skeleton was discovered in 1984 and is virtually complete, including a fragmentary skull, making Amargasaurus one of the best-known sauropods of its epoch. Amargasaurus was first described in 1991 and contains a single known species, Amargasaurus cazaui. It was a large animal, but small for a sauropod, reaching in length. Most distinctively, it sported two parallel rows of tall spines down its neck and back, taller than in any other known sauropod. In life, these spines could have stuck out of the body as solitary structures that supported a keratinous sheath. An alternate hypothesis, now more favored, postulates that they could have formed a scaffold supporting a skin sail. They might have been used for display, combat, or defense.
Amargasaurus was discovered in sedimentary rocks of the La Amarga Formation, which dates back to the Barremian and Aptian stages of the Early Cretaceous. A herbivore, it shared its environment with at least three other sauropod genera, which might have exploited different food sources in order to reduce competition. Amargasaurus probably fed at mid-height, as shown by the orientation of its inner ear and the articulation of its neck vertebrae, which suggest a habitual position of the snout above the ground and a maximum height of . Within the Sauropoda, Amargasaurus is classified as a member of the family Dicraeosauridae, which differs from other sauropods in showing shorter necks and smaller body sizes.
Description
Amargasaurus was small for a sauropod, measuring in length and weighing approximately . It followed the typical sauropod body plan, with a long tail and neck, a small head, and a barrel-shaped trunk supported by four column-like legs. The neck of Amargasaurus was shorter than in most other sauropods, a common trait within the Dicraeosauridae. Measuring in length, the neck corresponded to 136% of the length of the dorsal vertebral column. This is comparable to Dicraeosaurus (123%) but greater than in the extremely short-necked form Brachytrachelopan (75%). The neck consisted of thirteen cervical vertebrae, which were opisthocoelous (convex at the front and hollow at the back), forming ball-and-socket joints with neighboring vertebrae. The trunk was made out of nine dorsal and probably five fused sacral vertebrae. The foremost dorsals were opisthocoelous, while the remaining dorsals were amphiplatyan (flat on both ends). Robust transverse processes (lateral projections connecting to the ribs) indicate a strongly developed rib cage. The dorsal vertebrae of Amargasaurus and other dicraeosaurids lacked pleurocoels, the deep lateral excavations that were characteristic for other sauropods.
The most striking features of the skeleton were the extremely tall, upwardly projecting neural spines on the neck and anterior dorsal vertebrae. From the third cervical onward, these neural spines were bifurcated along their entire length, forming a double row. They were circular in cross section and tapered towards their tips. The tallest spines could be found on the middle part of the neck, where they reached on the 8th cervical. On the neck, they were bowed backward, projecting above the adjacent vertebra. A similar elongated neural spine has been described from the neck region of the closely related Bajadasaurus in 2019. Unlike in Amargasaurus, this spine was bowed frontward and broadened toward the tips. The last two dorsal vertebrae, the hip, and the foremost tail in Amargasaurus also had elongated spines; these were not bifurcated but flared into a paddle-shaped upper end. The pelvic region was relatively wide, judging from the long, laterally projecting transverse processes of the sacral vertebrae. The forelimbs were somewhat shorter than the hind limbs, as in related sauropods. Most of the hand and foot bones were not preserved, but Amargasaurus probably possessed five digits each as in all sauropods.
Only the rear part of the skull is preserved. It likely showed a horselike, broad snout equipped with pencil-like teeth, as seen in related sauropods for which more complete skulls are known. As in other dicraeosaurids, the (nostril opening) was situated in the posterior half of the skull, diagonally above the (eye opening), which was proportionally large. As in most other dinosaurs, the skull featured three additional openings (fenestrae). The , located below the orbit, was long and narrow. Behind the orbit was the , which in dicraeosaurids was uniquely small and can be seen when the skull is viewed from the side. This contrasts with other diapsid reptiles, where these openings were directed upward, thus being visible only in top view. The antorbital fenestra would have been located in front of the eye opening, although this region is not preserved. An unusual feature were small openings seen on the backside of the skull, the so-called parietal openings or fontanelles. In other tetrapods, these openings are usually seen only in juveniles and would close as the individual grows. Skull features shared with Dicraeosaurus but absent in most other sauropods included the fused and the notably long , bony extensions connecting the with the .
Discovery
The only known skeleton (specimen number MACN-N 15) was discovered in February 1984 by Guillermo Rougier during an expedition led by Argentine paleontologist José Bonaparte. This was the eighth expedition of the project "Jurassic and Cretaceous Terrestrial Vertebrates of South America", which was supported by the National Geographic Society and initiated in 1975 to improve on the sparse knowledge of the Jurassic and Cretaceous tetrapod life of South America. The same excursion uncovered the nearly complete skeleton of the horned theropod Carnotaurus. The discovery site is located in the La Amarga arroyo in the Picún Leufú Department of Neuquén Province in northern Patagonia, south of Zapala. The skeleton stems from sedimentary rocks of the La Amarga Formation, which dates to the Barremian through early Aptian stages of the Early Cretaceous, or around 130 to 120 million years ago.
The skeleton is reasonably complete and includes a partial skull. Sauropod skull bones are rarely found, and the Amargasaurus skull is only the second skull known from a member of the Dicraeosauridae. Major parts of the skeleton were found in their original anatomical position: the vertebral column of the neck and back, which consisted of 22 articulated vertebrae, was found connected to both the skull and the sacrum. Of the skull, only the temporal region and the braincase are preserved. The sacrum, despite being partly eroded prior to burial, is fairly complete. Most of the tail is missing, with three anterior, three middle, and one posterior vertebrae being preserved, along with fragments of several others. The shoulder girdle is known from the scapula (shoulder blade) and coracoid (which sits on the lower end of the scapula), while the pelvis is merely known from the ilium (the uppermost of the three pelvic bones). The limbs are equally fragmentary, with the manus (hand) and most of the pes (foot) missing. The skeleton is currently stored in the collection of the Bernardino Rivadavia Natural Sciences Museum in Buenos Aires.
The first, unofficial, mention of Amargasaurus as a new genus of dinosaur was published by Bonaparte in the 1984 Italian book Sulle Orme dei Dinosauri. Here, the species was designated as Amargasaurus groeberi, honoring Pablo Groeber, which was changed into Amargasaurus cazaui in the official description published several years later. The official description, written in Spanish, was published in 1991 by Leonardo Salgado and Bonaparte in the Argentinian scientific journal Ameghiniana. The name Amargasaurus alludes to the site of discovery, the La Amarga Arroyo. La Amarga is also the name of a nearby town, as well as the geologic formation the remains were recovered from. The word amarga is Spanish for "bitter", while sauros is Greek for "lizard". The one species (A. cazaui) is named in honor of Luis Cazau, a geologist with the YPF oil company, which at the time was state-owned. In 1983, Cazau informed Bonaparte's team about the paleontological significance of the La Amarga Formation, leading to the discovery of the skeleton. One year later, Salgado and Jorge O. Calvo published a second paper focusing on the description of the skull.
Classification
Amargasaurus is classified as a member of the Dicraeosauridae, a family ranked clade within the Diplodocoidea. Currently, this clade consists of nine species belonging to eight genera. These include Lingwulong shenqi from the Early or Middle Jurassic of China and four species from the Late Jurassic: Brachytrachelopan mesai from Argentina; Suuwassea emilieae from the Morrison Formation of the United States; and Dicraeosaurus hansemanni and Dicraeosaurus sattleri from the Tendaguru beds of Tanzania. Amargasaurus was the first dicraeosaurid known from the Cretaceous, although additional dicraeosaurids from the Lower Cretaceous have been described more recently, including Pilmatueia faundezi, Amargatitanis macni, and Bajadasaurus pronuspinax, which are all from Argentina. An unnamed specimen from Brazil indicates that this group persisted at least until the end of the Early Cretaceous. Most analyses find Dicraeosaurus and Brachytrachelopan to be more closely related to each other than to Amargasaurus.
Suuwassea is generally recovered as the most basal member of the family. A 2015 analysis by Tschopp and colleagues came to the preliminary result that two poorly known genera from the Morrison Formation, Dyslocosaurus polyonychius and Dystrophaeus viaemalae, might be additional members of the Dicraeosauridae.
Together with the Diplodocidae and the Rebbachisauridae, the Dicraeosauridae is nested inside the Diplodocoidea. All members of the Diplodocoidea are characterized by their box-shaped snout and narrow teeth restricted to the foremost portion of the jaws. Both the Dicraeosauridae and the Diplodocidae are characterized by bifurcated neural spines of the cervical and dorsal vertebra. In the Dicraeosauridae, the bifurcated neural spines were strongly elongated, a trend reaching its extreme in Amargasaurus.
The following cladogram by Gallina and colleagues (2019) shows the presumed relationships between members of the Dicraeosauridae:
Paleobiology
Vertebral spines
Both the function and the appearance in life of the extremely elongated and bifurcated vertebral spines remain elusive. Salgado and Bonaparte, in 1991, suggested the spines represented defense weapons against predators, arguing that they tapered towards their tips. They also could have served for display, perhaps for courtship or to intimidate rivals. Some subsequently published life restorations showed the double row of spines supporting two parallel skin sails. Gregory Paul, in 1994, considered this possibility unlikely, noting that neck sails would have reduced neck flexion, and that the spines were circular in cross-section rather than flattened as is the case in sail-bearing animals. Instead, he found that this shape indicates that the spines supported a keratinous sheath that would have extended the length of the spines in life. The spines could have been used for display or as weapons both against predators and members of the same species, as the animal might have been able to point its most anterior spines forward by bending its neck. He also hypothesized that the spines could have been clattered together to generate sound. Keratinous sheaths covering the spines were also shown in a 1999 skeletal restoration published by Salgado.
Jack Bailey, in 1997, argued that the spines resembled those of sail-bearing pelycosaurs like Dimetrodon. According to this author, Amargasaurus might also have possessed such a sail, which might have been used for display. Unlike those of pelycosaurs, the neural spines of Amargasaurus were bifurcated, forming a double row along the neck and back. As the space between both rows was merely , the existence of two parallel sails seems unlikely. Instead, Bailey suggested the spines represented a scaffold which was completely enveloped by a single skin. Neural spines from the penultimate dorsal vertebra to the foremost tail vertebrae also were strongly elongated, but different in structure, forming a single row of paddle-shaped projections. According to Bailey, these projections resembled those of modern humped ungulates such as the bison, indicating the presence of a fleshy hump above the hips. Bailey suggested similar humps for other dinosaurs with strongly elongated neural spines, such as Spinosaurus and Ouranosaurus.
Daniela Schwarz and colleagues, in 2007, concluded that the bifurcated neural spines of diplodocids and dicraeosaurids enclosed an air sac, which would have been connected to the lungs as part of the respiratory system. In Dicraeosaurus, this air sac (the so-called supravertebral diverticulum) would have rested on top of the neural arch and filled the entire space between the spines. In Amargasaurus the upper two-thirds of the spines would have been covered by a sheath of keratin, restricting the air sac to the space between the lower one-third of the spines. A cover of either keratin or skin is indicated striations on the surface of the spines similar to those of bony horn cores of today's bovids. In 2016, Mark Hallett and Mathew Wedel suggested that the backwards-directed spines might have been able to skewer predators when the neck was abruptly drawn backwards during an attack. A similar defense strategy is found in today's giant sable antelope and Arabian oryx, which can use their long, backwards directed horns to stab attacking lions. Apart from the possible function in defense, the spines may have been used for display, either for the intimidation of rivaling individuals or for courtship. Hallett and Wedel also hypothesized that rival males might have interlocked their spines for neck wrestling. Pablo Gallina and colleagues (2019) described the closely related Bajadasaurus, which had neural spines similar to those of Amargasaurus, and suggested that both genera employed them for defense. A defense function would have been especially effective in Bajadasaurus as the spines were directed forwards and would have reached past the tip of the snout, deterring predators. The keratinous sheath that likely covered the spines might have extended their length by 50%, as seen in some modern even-toed ungulates. Such extended sheaths would have made the delicate spines more resistant to damage—likely a critical threat, as the bases of the spines form the roof of the spinal cord.
In 2022, a detailed study by Ignacio A. Cerda and colleagues analyzing the structure, morphology, and microanatomy of the vertebral spines of Amargasaurus and an indeterminate dicraeosaurid (also from the La Amarga Formation) suggested that the spines were not covered in a keratinous sheath as previously believed. Osteohistology of the spines suggests that they were likely, if not exclusively, covered in a sail of skin. The spines are also highly vascularized and bear cyclical growth marks, adding credence to this theory.
Senses and posture
Paulina Carabajal and colleagues, in 2014, CT-scanned the skull, allowing for the generation of three-dimensional models of both the cranial endocast (the cast of the brain cavity) and the inner ear. Using these models, the cranial endocast was shown to encompass in volume. The inner ear was tall and wide. The lagena, the part containing the hair cells for hearing, was rather short, indicating that the sense of hearing would have been poorer in Amargasaurus than in other sauropods for which inner ears have been studied.
The first skeletal reconstructions show the skull in a near-horizontal posture. Salgado, in 1999, argued that such a posture would have been anatomically impossible due to the elongated neural spines of the neck vertebrae. Instead, he envisaged the head in a nearly vertical orientation. The habitual orientation of the head is usually reflected by the orientation of the semicircular canals of the inner ear, which housed the sense of balance (vestibular system). Using their three-dimensional model of the inner ear, Carabajal and colleagues suggested that the snout faced downwards at an angle of roughly 65° relative to the horizontal. A similar value has recently been proposed for the related Diplodocus. The neutral posture of the neck can be approximated based on how the cervical vertebrae attached to each other. According to Carabajal and colleagues, the neck was gently sloping downwards, so that the snout would have rested above the ground in neutral posture. In reality, neck posture would have varied according to the respective activities of the animals. Raising of the neck, e.g. for reaching an alert position, would have been constricted by the elongated neural spines, not permitting heights greater than .
Locomotion
Amargasaurus was quadrupedal (moved on four legs), and probably was unable to rear on its hind legs. Salgado and Bonaparte, in 1991, suggested that Amargasaurus was a slow walker, as both the forearms and lower legs were proportionally short, as a feature common to slow-moving animals. This was contradicted by Gerardo Mazzetta and Richard Fariña in 1999, who argued that Amargasaurus was capable of rapid locomotion. During locomotion, leg bones are strongly affected by bending moments, representing a limiting factor for the maximum speed of an animal. The leg bones of Amargasaurus were even more sturdy than those of today's white rhinoceros, which is adapted to galloping.
Life history
In a 2021 study, Guillermo Windholz and Ignacio Cerda obtained thin sections of the humerus, femur, and a rib of the Amargasaurus specimen in order to determine the number and spacing of lines of arrested growth (analogous to the growth rings of a tree). The rib showed the most complete record of lines of arrested growth, indicating that the Amargasaurus holotype individual was at least ten years old. In sauropodomorps, sexual maturity occurred long before adult size was reached. In the outer cortex (the most external layer of the bone when seen in cross section) of the Amargasaurus individual, lines of arrested growth are more abundant, indicating sexual maturity. However, an external fundamental system (a layer containing very closely spaced lines of arrested growth) is missing, possibly indicating that the individual was not yet fully grown, although it cannot be excluded that the external fundamental system was originally present but has since been eroded.
Paleoecology
Amargasaurus stems from sedimentary rocks of the La Amarga Formation, which is part of the Neuquén Basin and dates to the Barremian and late Aptian of the Early Cretaceous. Most vertebrate fossils, including Amargasaurus, have been found in the lowermost (oldest) part of the formation, the Puesto Antigual Member. This member is approximately in thickness and mainly composed of sandstones deposited by braided rivers. The Amargasaurus skeleton itself was recovered from a layer composed of sandy conglomerates. The sauropod fauna of the La Amarga Formation was diverse and included the basal rebbachisaurid Zapalasaurus, the dicraeosaurid Amargatitanis, and unnamed remains of basal titanosauriforms. The high diversity suggests that different sauropod species exploited different food sources in order to reduce competition. Basal Titanosauriforms showed proportionally longer necks, longer forelimbs, and broader tooth crowns than Dicraeosaurids and Rebbachisaurids, suggesting greater feeding heights. Amargasaurus probably fed above ground level at heights of up to , as evidenced by the anatomy of its neck and inner ear. Rebbachisaurids like Zapalasaurus presumably fed at ground-level, while basal Titanosauriforms exploited food sources at higher levels.
Other dinosaurs of the La Amarga Formation include an indeterminate stegosaur; predatory dinosaurs include the small ceratosaur Ligabueino, and the presence of a large tetanuran is indicated by teeth. Other than dinosaurs, the formation is notable for the cladotherian mammal Vincelestes, the only mammal known from the Early Cretaceous of South America. Crocodylomorphs are represented by the trematochampsid Amargasuchus – the holotype of this genus was found in association with the Amargasaurus bones.
| Biology and health sciences | Sauropods | Animals |
3535484 | https://en.wikipedia.org/wiki/Saffron%20%28color%29 | Saffron (color) | Saffron is a shade of yellow or orange, the colour of the tip of the saffron crocus thread, from which the spice saffron is derived. The hue of the spice saffron is primarily due to the carotenoid chemical crocin.
Etymology
The word saffron ultimately derives (via Arabic) from the Middle Iranian ja'far-. The name was used for the saffron spice in Middle English from c. 1200. As a colour name, it dates to the late 14th century.
Deep saffron approximates the colour of India saffron (also known as bhagwa or kesari).
In Rajasthani, this colour is called kesariya. The word derives its name from kesar, the Hindustani name for saffron, an important crop in Kashmir.
Religion
The color Saffron () is considered as a sacred color in Hinduism. According to Hindu mythology, Saffron (or Kesariya) is the color of Sunset (Sandhya) and Fire (Agni) which symbolises sacrifice, light, and quest of salvation. The color is worn by Hindu saints and ascetics as their devotion toward the religion. Many Hindu kingdoms and dynasties had Saffron color in their flag denoting the Sanātana Dharma, including Maratha Empire.
Hinduism, Jainism and Buddhism associate saffron with the pious renunciation of material life.
Buddhist monks in the Theravada tradition typically wear saffron robes (although occasionally maroon — the color normally worn by Vajrayana Buddhist monks — is worn). The tone of saffron typically worn by Theravada Buddhist monks is the lighter tone of saffron shown above.
Saffron holds symbolic meaning in Sikhism, representing spirit and sacrifice. Originally a shade of yellow called basanti, the field of the modern Nishan Sahib is saffron. Turbans worn by Sikhs most often are blue or white, but basanti colour is common.
Political & religious uses
In politics, it was used by the Indian independence movement, and it was chosen as one of the three colours of the Indian national flag after independence in 1947, and is used by Hindus. India saffron, representing courage and sacrifice, was chosen for one of the three bands of the National Flag of India, along with white (peace and truth) and what is now called India green (faith and chivalry). The Flag of India is officially described in the Flag Code of India as follows:The colour of the top panel shall be India saffron (Kesari) and that of the bottom panel shall be India green. The middle panel shall be white, bearing at its centre the design of Ashoka Chakra in navy blue colour with 24 equally spaced spokes.Sarvepalli Radhakrishnan, who later became India's first Vice President and second President, described the significance of the Indian National Flag as follows:
The use of saffron in the national flag and as political symbolism has been opposed. One line of opposition asserts that the color is sacred and should not be politicized. Another source of opposition comes from Islamists who claim the color is forbidden in Islam and strongly prohibited to be worn by the males.
Basanti turbans are associated with the Khalistan movement in the Punjab region of Pakistan and India. Even otherwise Basanti turbans are commonly used by Sikhs and not all wearing Basanti turban should be associated with separatist movement.
Because Therevada Buddhist monks were at the forefront of the 2007 Burmese anti-government protests, the uprising has been referred to as the Saffron Revolution by some in the international media.
Hindu nationalism
The saffron flag (bhagwa dhwaj) of the medieval Hindu warrior Chhatrapati Shivaji Maharaj was held in high esteem by the Hindu Mahasabha and Rashtriya Swayamsevak Sangh (RSS) in the 1920s as a representative of Hindu resurgence and militaristic tradition. The saffron flag was the "true guru" to which Hedgewar demanded obeisance from the RSS members. "The Gerva [saffron] Flag shall be the flag of the Hindu nation. With its Om, the Swastik and the Sword, it appeals to the sentiments cherished by our race since the Vaidik [Vedic] days," he said.
The Bharatiya Jana Sangh and its successor Bharatiya Janata Party (BJP) both used saffron as their colour. The BJP used a saffron lotus on its flag, along with a green side band that possibly reflected accommodation with Islam. The Vishva Hindu Parishad (VHP), a Hindu religious body affiliated to the RSS, also used saffron as its predominant colour, with its ascetic leaders clad in saffron robes and the lay leaders wearing saffron scarves. During the Ram Janmabhoomi movement in the 1990s, the VHP and its affiliate Bajrang Dal distributed saffron flags and saffron headbands to their followers by the millions.
The predominance of the saffron symbolism in the BJP and its allies led to the BJP being referred to as the 'saffron party' in the 1990s, and the term 'saffronisation' came to be used describe the increasing influence of Hindu nationalism in party politics. This period saw phrases such as the "saffronisation of the coastal belt", "saffronisation of Karnataka" and "saffronisation of the Congress(I)". Academic and non-academic scholars wrote books with titles involving 'saffron' to refer to Hindu nationalism: Brotherhood in Saffron, Khaki Shorts and Saffron Flags, The Saffron Wave, and The Saffron Swastika.
Clothing
Saffron-coloured cloth had a history of use among the Gaelic-Irish. A saffron kilt is worn by the pipers of certain Irish regiments in the British Army, and the saffron léine in the defence forces of the Republic of Ireland. The latter garment is also worn by some Irish and Irish-American men as an item of national costume (though most wear kilts, believing them to be Irish). Its colour varies from a true saffron orange to a range of dull mustard and yellowish-brown hues.
The Antrim GAA teams are nicknamed "The Saffrons" because of the saffron-coloured kit which they play in. The Old Irish word for saffron, cróc, derives directly from the Latin Crocus sativus. In Ireland between the 14th and 17th centuries, men wore léinte (singular léine), loose saffron-coloured shirts that reached down to mid-thigh or the knee. (see Irish clothing).
Literature
The colour saffron is associated with the goddess of dawn (Eos in Greek mythology and Aurora in Roman mythology) in classical literature:
Homer's Iliad:Now when Dawn in robe of saffron was hastening from the streams of Okeanos, to bring light to mortals and immortals, Thetis reached the ships with the armor that the god had given her. (19.1)Virgil's Aeneid:Aurora now had left her saffron bed,
And beams of early light the heav'ns o'erspread,
When, from a tow'r, the queen, with wakeful eyes,
Saw day point upward from the rosy skies.
Other media
The lyrics of Donovan's 1966 song, "Mellow Yellow" repeat the line, "I'm just mad about Saffron".
In the Pokémon franchise, in the region of Kanto there is a city named Saffron City.
The Gates is a site-specific art installation by Christo and Jeanne-Claude. The artists installed 7,503 metal "gates" along 23 miles (37 km) of pathways in Central Park in New York City. From each gate hung a flag-shaped piece of deep saffron-coloured nylon fabric. The exhibit ran from February 12, 2005, through February 27, 2005.
Saffron Monsoon is a character in Absolutely Fabulous.
In nature
Plants
Byzantine meadow saffron (Colchicum × byzantinum) is a hybrid flowering plant.
Cape saffron (Cassine peragua) is a flowering tree with saffron-coloured bark.
Cobra saffron (Mesua ferrea) is a tree found in southern Asia.
Meadow saffron (Colchicum autumnale) is a flowering plant found in Europe.
Mediterranean meadow saffron (Colchicum cupanii) is a flowering plant found in central Mediterranean basin.
Saffron buckwheat (Eriogonum crocatum) is a species of wild buckwheat endemic to the Conejo Valley.
Saffron spice is derived from the flowers of the plant named saffron crocus (Crocus sativus).
Saffron plum (Sideroxylon celastrinum) is a flowering plant found in North, Central, and South America.
Saffron thistle (Carthamus lanatus) is a thistle native to the Mediterranean basin.
Spring meadow saffron (Colchicum bulbocodium) is a flowering alpine plant found in Europe.
Steven's meadow saffron (Colchicum stevenii) is a flowering plant found in the eastern Mediterranean.
Birds
The saffron-billed sparrow (Arremon flavirostris) is a bird found in South America.
The saffron-breasted prinia (Prinia hypoxantha) is a passerine bird found in eastern South Africa and Swaziland.
The saffron-crested tyrant-manakin (Neopelma chrysocephalum) is a bird found in the Guianas, southern Venezuela, and the northwestern Amazon basin.
The saffron-crowned tanager (Tangara xanthocephala) is a bird found in the montane forests of South America.
The saffron-headed parrot (Pyrilia pyrilia) is a parrot found in the montane forests of South America.
The saffron finch (Sicalis flaveola) is a tanager from South America, and is common in both open and semi-open areas in lowlands outside the Amazon basin.
The saffron siskin (Spinus siemiradzkii) is a finch found in Ecuador and Peru.
The saffron toucanet (Pteroglossus bailloni) is a toucan from South American's Atlantic Forest.
Aquatic animals
The saffron cod (Eleginus gracilis) is a commercially harvested fish in the North Pacific.
The saffron-coloured clam (Tridacna crocea) is a bivalve found in the Indo-Pacific region.
The saffron shiner (Notropis rubricroceus) is a fish found in Tennessee River drainage.
Amphibians
The saffron-bellied frog (Chaperina fusca) is a frog found in the Malay Peninsula, Borneo, and the Philippines.
Insects
The saffron-winged meadowhawk (Sympetrum costiferum) is a dragonfly found in North America.
The saffron beetle (Calosoma schayeri) is a beetle found in Australia.
The saffron sapphire (Iolaus pallene) is a butterfly found in Africa.
The saffron skipper (Poanes aaroni) is a skipper found in North America.
Fungi
False saffron milkcap (Lactarius deterrimus) is a fungus found in Europe and Asia.
Saffron milk cap (Lactarius deliciosus) is an edible fungus found in Europe.
Saffron ringless amanita (Amanita crocea) is a Amantia found in Europe.
Viruses
The Saffron Scourge is another name for yellow fever.
| Physical sciences | Colors | Physics |
1244402 | https://en.wikipedia.org/wiki/Gas%20balloon | Gas balloon | A gas balloon is a balloon that rises and floats in the air because it is filled with a gas lighter than air (such as helium or hydrogen). When not in flight, it is tethered to prevent it from flying away and is sealed at the bottom to prevent the escape of gas. A gas balloon may also be called a Charlière for its inventor, the Frenchman Jacques Charles. Today, familiar gas balloons include large blimps and small latex party balloons. For nearly 200 years, well into the 20th century, manned balloon flight utilized gas balloons before hot-air balloons became dominant. Without power, heat or fuel, untethered flights of gas balloons depended on the skill of the pilot. Gas balloons have greater lift for a given volume, so they do not need to be so large, and they can stay up for much longer than hot air balloons.
History
The first gas balloon made its flight in August 1783. Designed by professor Jacques Charles and Les Frères Robert, it carried no passengers or cargo. On 1 December 1783, their second hydrogen-filled balloon made a manned flight piloted by Jacques Charles and Nicolas-Louis Robert. This occurred ten days after the first manned flight in a Montgolfier hot air balloon.
The next project of Jacques Charles and the Robert brothers was La Caroline, an elongated steerable craft that followed Jean Baptiste Meusnier's proposals for a dirigible balloon, incorporating internal ballonnets (air cells), a rudder and a method of propulsion. On September 19, 1784 the brothers and M. Collin-Hullin flew for 6 hours 40 minutes, covering 186 km from Paris to Beuvry near Béthune. This was the first flight over 100 km.
Gas balloons remained popular throughout the age before powered flight. Filled with hydrogen or coal gas, they were able to fly higher, further and more economically than hot-air balloons. The altitude was controlled with ballast weights that were dropped if the balloon got too low; in order to land some lifting gas was vented through a valve. Tethered manned gas balloons were used for observation purposes in the Napoleonic Wars (to very limited extent), in the American Civil War (flown by Thaddeus Lowe) and in World War I by aviators wearing parachutes. Throughout the 19th century, they were popular as objects of public fascination among hobbyists and show performers, such as the Blanchards.
Throughout the mid 20th century, spherical free gas balloons were used by the United States Navy to train airship crews.
Gas ballooning has been popular in Europe, most notably in Germany, using hydrogen as a lifting gas. Gas balloon clubs exist throughout the country. Rough estimates show 150 active gas pilots in Europe. In contrast, gas ballooning in the USA might have at most 30 active pilots who typically fly only once a year at the Albuquerque International Balloon Fiesta in October.
Aerophile is the world's largest lighter-than-air carrier, flying 300,000 passengers every year through its eight tethered gas balloon operations in Walt Disney World, San Diego Zoo Safari Park, Smoky Mountains & Irvine in the US and Paris, Disneyland Paris and Parc du Petit Prince in France.
Records
On October 24, 2014, Alan Eustace, a former Google executive, made a jump from the stratosphere, breaking Felix Baumgartner's 2012 world record. The launch-point for his jump was from an abandoned runway in Roswell, New Mexico, where he began his balloon-powered ascent early that morning. He reached a reported maximum altitude of 135,908 feet (41.425 km; 25.7402 mi), but the final number submitted to the World Air Sports Federation was 135,889.108 feet (41.419000 km; 25.7365735 mi). The balloon used for the feat was manufactured by the Balloon Facility of the Tata Institute of Fundamental Research, Hyderabad, India. Eustace in his pressure suit hung tethered under the balloon, without the kind of capsule used by Felix Baumgartner. Eustace started his fall by using an explosive device to separate from the helium balloon.
The previous altitude record for a manned balloon flight was set at 39.045 kilometers on October 14, 2012 by Felix Baumgartner breaking a record of 34.7 kilometers on May 4, 1961 by Malcolm Ross and Victor Prather in a balloon launched from the deck of the in the Gulf of Mexico.
The altitude record for an unmanned balloon is 53.7 kilometers. It was reached by a stratospheric balloon manufactured by JAXA with a volume of 80,000 m³, launched in September 2013 from in Hokkaido, Japan. This is the greatest height ever obtained by an atmospheric vehicle. Only rockets, rocket planes, and ballistic projectiles have flown higher.
In 2015, pilots Leonid Tiukhtyaev and Troy Bradley arrived safely in Baja California, Mexico, after a journey of 10,711 km. The two men, originally from Russia and the United States of America respectively, started in Japan and flew with a helium balloon over the Pacific. In 160 hours and 34 minutes, the balloon, named "Two Eagles", arrived in Mexico, setting the longest distance and duration records for gas balloons.
On other planets
The Soviet Union space probes Vega 1 and Vega 2 each dropped a helium balloon with scientific experiments into the atmosphere of Venus in 1985. The balloons first entered the atmosphere and descended to about 50 km, then inflated for level flight. Otherwise the flight was uncontrolled. Each balloon relayed wind and atmospheric conditions for 46 hours of a possible 60-hour electric battery power budget.
| Technology | Types of aircraft | null |
1244926 | https://en.wikipedia.org/wiki/Graveyard%20orbit | Graveyard orbit | A graveyard orbit, also called a junk orbit or disposal orbit, is an orbit that lies away from common operational orbits. One significant graveyard orbit is a supersynchronous orbit well beyond geosynchronous orbit. Some satellites are moved into such orbits at the end of their operational life to reduce the probability of colliding with operational spacecraft and generating space debris.
Overview
A graveyard orbit is used when the change in velocity required to perform a de-orbit maneuver is too large. De-orbiting a geostationary satellite requires a delta-v of about , whereas re-orbiting it to a graveyard orbit only requires about .
For satellites in geostationary orbit and geosynchronous orbits, the graveyard orbit is a few hundred kilometers beyond the operational orbit. The transfer to a graveyard orbit beyond geostationary orbit requires the same amount of fuel as a satellite needs for about three months of stationkeeping. It also requires a reliable attitude control during the transfer maneuver. While most satellite operators plan to perform such a maneuver at the end of their satellites' operational lives, through 2005 only about one-third succeeded. Given the economic value of the positions at geosynchronous altitude, unless premature spacecraft failure precludes it, satellites are moved to a graveyard orbit prior to decommissioning.
According to the Inter-Agency Space Debris Coordination Committee (IADC) the minimum perigee altitude beyond the geostationary orbit is:
where is the solar radiation pressure coefficient and is the aspect area [m2] to mass [kg] ratio of the satellite. This formula includes about 200 km for the GEO-protected zone to also permit orbit maneuvers in GEO without interference with the graveyard orbit. Another of tolerance must be allowed for the effects of gravitational perturbations (primarily solar and lunar). The remaining part of the equation considers the effects of the solar radiation pressure, which depends on the physical parameters of the satellite.
In order to obtain a license to provide telecommunications services in the United States, the Federal Communications Commission (FCC) requires all geostationary satellites launched after March 18, 2002, to commit to moving to a graveyard orbit at the end of their operational lives. U.S. government regulations require a boost, , of about . In 2023 DISH received the first-ever fine by the FCC for failing to de-orbit its EchoStar VII satellite according to the terms of its license.
A spacecraft moved to a graveyard orbit will typically be passivated.
Uncontrolled objects in a near geostationary [Earth] orbit (GEO) exhibit a 53-year cycle of orbital inclination due to the interaction of the Earth's tilt with the lunar orbit. The orbital inclination varies ± 7.4°, at up to 0.8°pa.
Disposal orbit
While the standard geosynchronous satellite graveyard orbit results in an expected orbital lifetime of millions of years, the increasing number of satellites, the launch of microsatellites, and the FCC approval of large megaconstellations of thousands of satellites for launch by 2022 necessitates new approaches for deorbiting to assure earlier removal of the objects once they have reached end-of-life. Contrary to GEO graveyard orbits requiring three months' worth of fuel (delta-V of 11 m/s) to reach, large satellite networks in LEO require orbits that passively decay into the Earth's atmosphere. For example, both OneWeb and SpaceX have committed to the FCC regulatory authorities that decommissioned satellites will decay to a lower orbita disposal orbitwhere the satellite orbital altitude would decay due to atmospheric drag and then naturally reenter the atmosphere and burn up within one year of end-of-life.
| Physical sciences | Orbital mechanics | Astronomy |
1245859 | https://en.wikipedia.org/wiki/Fayalite | Fayalite | Fayalite (, commonly abbreviated to Fa) is the iron-rich end-member of the olivine solid-solution series. In common with all minerals in the olivine group, fayalite crystallizes in the orthorhombic system (space group Pbnm) with cell parameters a 4.82 Å, b 10.48 Å and c 6.09 Å.
Fayalite forms solid solution series with the magnesium olivine endmember forsterite (Mg2SiO4) and also with the manganese rich olivine endmember tephroite (Mn2SiO4).
Iron rich olivine is a relatively common constituent of acidic and alkaline igneous rocks such as volcanic obsidians, rhyolites, trachytes and phonolites and plutonic quartz syenites where it is associated with amphiboles. Its main occurrence is in ultramafic volcanic and plutonic rocks and less commonly in felsic plutonic rocks and rarely in granite pegmatite. It also occurs in lithophysae in obsidian. It also occurs in medium-grade thermally metamorphosed iron-rich sediments and in impure carbonate rocks.
Fayalite is stable with quartz at low pressures, whereas more magnesian olivine is not, because of the reaction olivine + quartz = orthopyroxene. Iron stabilizes the olivine + quartz pair. The pressure and compositional dependence of the reaction can be used to calculate constraints on pressures at which assemblages of olivine + quartz formed.
Fayalite can also react with oxygen to produce magnetite + quartz: the three minerals together make up the "FMQ" oxygen buffer. The reaction is used to control the fugacity of oxygen in laboratory experiments. It can also be used to calculate the fugacity of oxygen recorded by mineral assemblages in metamorphic and igneous processes.
At high pressure, fayalite undergoes a phase transition to ahrensite, the iron-bearing analogue of ringwoodite, i.e., contrary to forsterite there is no intermediate form analogous to wadsleyite; under the conditions prevailing in the upper mantle of the Earth, the transition would occur at ca. 6–7 GPa, i.e., at substantially lower pressure than the phase transitions of forsterite. In high-pressure experiments, the transformation may be delayed, so that it may remain stable to pressures of almost 35 GPa (see fig.), at which point it may become amorphous rather than take on a crystalline structure such as ahrensite.
The name fayalite is derived from Faial (Fayal) Island in the Azores where it was first described in 1840.
| Physical sciences | Silicate minerals | Earth science |
1246386 | https://en.wikipedia.org/wiki/Unbibium | Unbibium | Unbibium, also known as element 122 or eka-thorium, is a hypothetical chemical element; it has placeholder symbol Ubb and atomic number 122. Unbibium and Ubb are the temporary systematic IUPAC name and symbol respectively, which are used until the element is discovered, confirmed, and a permanent name is decided upon. In the periodic table of the elements, it is expected to follow unbiunium as the second element of the superactinides and the fourth element of the 8th period. Similarly to unbiunium, it is expected to fall within the range of the island of stability, potentially conferring additional stability on some isotopes, especially 306Ubb which is expected to have a magic number of neutrons (184).
Despite several attempts, unbibium has not yet been synthesized, nor have any naturally occurring isotopes been found to exist. There are currently no plans to attempt to synthesize unbibium. In 2008, it was claimed to have been discovered in natural thorium samples, but that claim has now been dismissed by recent repetitions of the experiment using more accurate techniques.
Chemically, unbibium is expected to show some resemblance to cerium and thorium. However, relativistic effects may cause some of its properties to differ; for example, it is expected to have a ground state electron configuration of [Og] 7d1 8s2 8p1 or [Og] 8s2 8p2, despite its predicted position in the g-block superactinide series.
Introduction
History
Synthesis attempts
Fusion-evaporation
Two attempts were made to synthesize unbibium in the 1970s, both propelled by early predictions on the island of stability at N = 184 and Z > 120, and in particular whether superheavy elements could potentially be naturally occurring. The first attempts to synthesize unbibium were performed in 1972 by Flerov et al. at the Joint Institute for Nuclear Research (JINR), using the heavy-ion induced hot fusion reactions:
+ → * → no atoms
Another unsuccessful attempt to synthesize unbibium was carried out in 1978 at the GSI Helmholtz Center, where a natural erbium target was bombarded with xenon-136 ions:
+ → * → no atoms
No atoms were detected and a yield limit of 5 nb (5,000 pb) was measured. Current results (see flerovium) have shown that the sensitivity of these experiments were too low by at least 3 orders of magnitude. In particular, the reaction between 170Er and 136Xe was expected to yield alpha emitters with half-lives of microseconds that would decay down to isotopes of flerovium with half-lives perhaps increasing up to several hours, as flerovium is predicted to lie near the center of the island of stability. After twelve hours of irradiation, nothing was found in this reaction. Following a similar unsuccessful attempt to synthesize unbiunium from 238U and 65Cu, it was concluded that half-lives of superheavy nuclei must be less than one microsecond or the cross sections are very small. More recent research into synthesis of superheavy elements suggests that both conclusions are true.
In 2000, the Gesellschaft für Schwerionenforschung (GSI) Helmholtz Center for Heavy Ion Research performed a very similar experiment with much higher sensitivity:
+ → * → no atoms
These results indicate that the synthesis of such heavier elements remains a significant challenge and further improvements of beam intensity and experimental efficiency is required. The sensitivity should be increased to 1 fb in the future for more quality results.
Compound nucleus fission
Several experiments studying the fission characteristics of various superheavy compound nuclei such as 306Ubb were performed between 2000 and 2004 at the Flerov Laboratory of Nuclear Reactions. Two nuclear reactions were used, namely 248Cm + 58Fe and 242Pu + 64Ni. The results reveal how superheavy nuclei fission predominantly by expelling closed shell nuclei such as 132Sn (Z = 50, N = 82). It was also found that the yield for the fusion-fission pathway was similar between 48Ca and 58Fe projectiles, suggesting a possible future use of 58Fe projectiles in superheavy element formation.
Claimed discovery as a naturally occurring element
In 2008, a group led by Israeli physicist Amnon Marinov at the Hebrew University of Jerusalem claimed to have found single atoms of unbibium-292 in naturally occurring thorium deposits at an abundance of between 10−11 and 10−12 relative to thorium. This was the first time in 69 years that a new element had been claimed to be discovered in nature, after Marguerite Perey's 1939 discovery of francium. The claim of Marinov et al. was criticized by the scientific community, and Marinov says he has submitted the article to the journals Nature and Nature Physics but both turned it down without sending it for peer review. The unbibium-292 atoms were claimed to be superdeformed or hyperdeformed isomers, with a half-life of at least 100 million years.
A criticism of the technique, previously used in purportedly identifying lighter thorium isotopes by mass spectrometry, was published in Physical Review C in 2008. A rebuttal by the Marinov group was published in Physical Review C after the published comment.
A repeat of the thorium experiment using the superior method of accelerator mass spectrometry (AMS) failed to confirm the results, despite a 100-fold better sensitivity. This result throws considerable doubt on the results of the Marinov collaboration with regards to their claims of long-lived isotopes of thorium, roentgenium, and unbibium. Current understanding of superheavy elements indicates that it is very unlikely for any traces of unbibium to persist in natural thorium samples.
Naming
Using Mendeleev's nomenclature for unnamed and undiscovered elements, unbibium should instead be known as eka-thorium. After the recommendations of the IUPAC in 1979, the element has since been largely referred to as unbibium with the atomic symbol of (Ubb), as its temporary name until the element is officially discovered and synthesized, and a permanent name is decided on. Scientists largely ignore this naming convention, and instead simply refer to unbibium as "element 122" with the symbol of (122), or sometimes even E122 or 122.
Prospects for future synthesis
Every element from mendelevium onward was produced in fusion-evaporation reactions, culminating in the discovery of the heaviest known element oganesson in 2002 and most recently tennessine in 2010. These reactions approached the limit of current technology; for example, the synthesis of tennessine required 22 milligrams of 249Bk and an intense 48Ca beam for six months. The intensity of beams in superheavy element research cannot exceed 1012 projectiles per second without damaging the target and detector, and producing larger quantities of increasingly rare and unstable actinide targets is impractical.
Consequently, future experiments must be done at facilities such as the superheavy element factory (SHE-factory) at the Joint Institute for Nuclear Research (JINR) or RIKEN, which will allow experiments to run for longer stretches of time with increased detection capabilities and enable otherwise inaccessible reactions.
It is possible that fusion-evaporation reactions will not be suitable for the discovery of unbibium or heavier elements. Various models predict increasingly short alpha and spontaneous fission half-lives for isotopes with Z = 122 and N ~ 180 on the order of microseconds or less, rendering detection nearly impossible with current equipment. The increasing dominance of spontaneous fission also may sever possible ties to known nuclei of livermorium or oganesson and make identification and confirmation more difficult; a similar problem occurred in the road to confirmation of the decay chain of 294Og which has no anchor to known nuclei. For these reasons, other methods of production may need to be researched such as multi-nucleon transfer reactions capable of populating longer-lived nuclei. A similar switch in experimental technique occurred when hot fusion using 48Ca projectiles was used instead of cold fusion (in which cross sections decrease rapidly with increasing atomic number) to populate elements with Z > 113.
Nevertheless, several fusion-evaporation reactions leading to unbibium have been proposed in addition to those already tried unsuccessfully, though no institution has immediate plans to make synthesis attempts, instead focusing first on elements 119, 120, and possibly 121. Because cross sections increase with asymmetry of the reaction, a chromium beam would be most favorable in combination with a californium target, especially if the predicted closed neutron shell at N = 184 could be reached in more neutron-rich products and confer additional stability. In particular, the reaction between and would generate the compound nucleus and reach the shell at N = 184, though the analogous reaction with a target is believed to be more feasible because of the presence of unwanted fission products from and difficulty in accumulating the required amount of target material. One possible synthesis of unbibium could occur as follows:
+ → + 3
Should this reaction be successful and alpha decay remain dominant over spontaneous fission, the resultant 300Ubb would decay through 296Ubn which may be populated in cross-bombardment between 249Cf and 50Ti. Although this reaction is one of the most promising options for the synthesis of unbibium in the near future, the maximum cross section is predicted to be 3 fb, one order of magnitude lower than the lowest measured cross section in a successful reaction. The more symmetrical reactions 244Pu + 64Ni and 248Cm + 58Fe have also been proposed and may produce more neutron-rich isotopes. With increasing atomic number, one must also be aware of decreasing fission barrier heights, resulting in lower survival probability of compound nuclei, especially above the predicted magic numbers at Z = 126 and N = 184.
Predicted properties
Nuclear stability and isotopes
The stability of nuclei decreases greatly with the increase in atomic number after plutonium, the heaviest primordial element, so that all isotopes with an atomic number above 101 decay radioactively with a half-life under a day. No elements with atomic numbers above 82 (after lead) have stable isotopes. Nevertheless, because of reasons not very well understood yet, there is a slight increased nuclear stability around atomic numbers 110–114, which leads to the appearance of what is known in nuclear physics as the "island of stability". This concept, proposed by University of California professor Glenn Seaborg, explains why superheavy elements last longer than predicted.
In this region of the periodic table, N = 184 has been suggested as a closed neutron shell, and various atomic numbers have been proposed as closed proton shells, such as Z = 114, 120, 122, 124, and 126. The island of stability would be characterized by longer half-lives of nuclei located near these magic numbers, though the extent of stabilizing effects is uncertain due to predictions of weakening of the proton shell closures and possible loss of double magicity. More recent research predicts the island of stability to instead be centered at beta-stable copernicium isotopes 291Cn and 293Cn, which would place unbibium well above the island and result in short half-lives regardless of shell effects. The increased stability of elements 112–118 has also been attributed to the oblate shape of such nuclei and resistance to spontaneous fission. The same model also proposes 306Ubb as the next spherical doubly magic nucleus, thus defining the true island of stability for spherical nuclei.
A quantum tunneling model predicts the alpha-decay half-lives of unbibium isotopes 284–322Ubb to be on the order of microseconds or less for all isotopes lighter than 315Ubb, highlighting a significant challenge in experimental observation of this element. This is consistent with many predictions, though the exact location of the 1 microsecond border varies by model. Additionally, spontaneous fission is expected to become a major decay mode in this region, with half-lives on the order of femtoseconds predicted for some even–even isotopes due to minimal hindrance resulting from nucleon pairing and loss of stabilizing effects farther away from magic numbers. A 2016 calculation on the half-lives and probable decay chains of isotopes 280–339Ubb yields corroborating results: 280–297Ubb will be proton unbound and possibly decay by proton emission, 298–314Ubb will have alpha half-lives on the order of microseconds, and those heavier than 314Ubb will predominantly decay by spontaneous fission with short half-lives. For the lighter alpha emitters that may be populated in fusion-evaporation reactions, some long decay chains leading down to known or reachable isotopes of lighter elements are predicted. Additionally, the isotopes 308–310Ubb are predicted to have half-lives under 1 microsecond, too short for detection as a result of significantly lower binding energy for neutron numbers immediately above the N = 184 shell closure. Alternatively, a second island of stability with total half-lives of approximately 1 second may exist around Z ~ 124 and N ~ 198, though these nuclei will be difficult or impossible to reach using current experimental techniques. However, these predictions are strongly dependent on the chosen nuclear mass models, and it is unknown which isotopes of unbibium will be most stable. Regardless, these nuclei will be hard to synthesize as no combination of obtainable target and projectile can provide enough neutrons in the compound nucleus. Even for nuclei reachable in fusion reactions, spontaneous fission and possibly also cluster decay might have significant branches, posing another hurdle to identification of superheavy elements as they are normally identified by their successive alpha decays.
Chemical
Unbibium is predicted to be similar in chemistry to cerium and thorium, which likewise have four valence electrons above a noble gas core, although it may be more reactive. Additionally, unbibium is predicted to belong to a new block of valence g-electron atoms, although the 5g orbital is not expected to start filling until about element 125. The predicted ground-state electron configuration of unbibium is either [Og] 7d1 8s2 8p1 or 8s2 8p2, in contrast to the expected [Og] 5g2 8s2 in which the 5g orbital starts filling at element 121. (The ds2p and s2p2 configurations are expected to be only separated by about 0.02 eV.) In the superactinides, relativistic effects might cause a breakdown of the Aufbau principle and create overlapping of the 5g, 6f, 7d and 8p orbitals; experiments on the chemistry of copernicium and flerovium provide strong indications of the increasing role of relativistic effects. As such, the chemistry of elements following unbibium becomes more difficult to predict.
Unbibium would most likely form a dioxide, UbbO2, and tetrahalides, such as UbbF4 and UbbCl4. The main oxidation state is predicted to be +4, similar to cerium and thorium. A first ionization energy of 5.651 eV and second ionization energy of 11.332 eV are predicted for unbibium; this and other calculated ionization energies are lower than the analogous values for thorium, suggesting that unbibium will be more reactive than thorium.
| Physical sciences | Periods | Chemistry |
1246694 | https://en.wikipedia.org/wiki/Chinese%20alligator | Chinese alligator | The Chinese alligator (; ), also known as the Yangtze alligator (), China alligator, or historically the muddy dragon, is a crocodilian endemic to China. It and the American alligator (A. mississippiensis) are the only living species in the genus Alligator of the family Alligatoridae. Dark gray or black in color with a fully armored body, the Chinese alligator grows to in length and weighs as an adult. It brumates in burrows in winter and is nocturnal in summer. Mating occurs in early summer, with females most commonly producing 20–30 eggs, which are smaller than those of any other crocodilian. The species is an opportunistic feeder, primarily eating fish and invertebrates. A vocal species, adults bellow during the mating season and young vocalize to communicate with their parents and other juveniles. Captive specimens have reached age 70, and wild specimens can live past 50.
Living in bodies of fresh water, the Chinese alligator's range is restricted to six regions in the province of Anhui, as well as possibly the provinces of Jiangsu and Zhejiang. Originally living as far away from its current range as Japan, the species previously had a wide range and population, but beginning in 6000 BC, multiple threats, such as habitat destruction, caused the species' population and range to decline. The population in the wild was about 1,000 in the 1970s, decreased to below 130 in 2001, and grew after 2003, with its population being about 300 as of 2017. Listed as critically endangered by the International Union for Conservation of Nature, multiple conservation actions have been taking place for this species.
The Chinese alligator has been a part of Chinese literature since the third century. In the late 13th century, Marco Polo became the first person outside of China to write about it. In some writings, the Chinese alligator has been associated with the Chinese dragon. Many pieces of evidence suggest that the Chinese alligator was an inspiration for the Chinese dragon.
History and taxonomy
The oldest definitive record of the Chinese alligator is from the late Pliocene of Japan, around 3 million years old. Pleistocene fossils show that its range was once much more extensive, extending northwards to Shandong and southwards to the Taiwan Strait.
Chinese alligators were mentioned in Chinese literature very early; for example, in the Classic of Poetry, whose poems were composed between the 11th and 7th centuries BCE. Marco Polo was the first person outside of China to write about the alligator, when he came to China and saw it in the late 1200s. He said that the alligator lived in "caverns" in the day and hunted at night, and that humans targeted its meat and skin, with its gall bladder having multiple medical purposes. He stated that it was found in lakes, rivers, and springs in the province "Karazan". In 1656, Martino Martini, a priest, wrote that the Chinese alligator lived in the river Yangtze and was "much feared by the local residents". Unlike Polo, Martini wrote his description using information from Chinese literature. Chinese alligators were later thought to give Buddhist priests merit if the priests were to buy alligators held in captivity and release them. In 1869, Robert Swinhoe saw a Chinese alligator in an exhibit in Shanghai and wrote the following year:
Classification
The Chinese alligator was scientifically described by French naturalist Albert-Auguste Fauvel in 1879 as Alligator sinensis; though Fauvel only noticed mentions of them in Chinese literature since about 222–227 CE. The genus Alligator had previously contained only the American alligator since its creation in 1807. Fauvel wrote a detailed description of the species in a book titled Alligators in China: Their History, Description & Identification, including information about its historical account. In 1947, it was suggested to group the Chinese alligator in a separate genus from its American relative, due to the Chinese alligator's bony plate on its upper eyelid. This bony plate is present in caimans, but is rarely present in the American alligator. At the time, the plate was thought to not appear in the American alligator at all. This produced the belief that the Chinese alligator's relationship with other crocodilians was between caimans and American alligators. Paulus Edward Pieris Deraniyagala described the genus Caigator the same year, which only contained the Chinese alligator, making its scientific name Caigator sinensis. However, paleontology has shown that the Chinese alligator has evolved from other now-extinct members of the genus Alligator. This and the fact that the American alligator does infrequently have a bony plate on its eyelid have caused Caigator sinensis to now be classified as a synonym of Alligator sinensis. There is still not a consensus among biologists that the American and Chinese alligators belong to the same genus, despite multiple studies comparing the biochemistry, histology, and various other aspects of the two crocodilians.
The evolutionary relationships of alligators can be shown in the cladogram below:
Etymology
The genus, Alligator, is based on the Spanish word . The specific name, sinensis, is from the Latin plural possessive , meaning "belonging to China".
Description
One of the smallest species of crocodilians, the Chinese alligator attains a length of and weight of as an adult. Females are roughly three-quarters the length of males. It is less than half the size of the American alligator, which typically grows to a length of for males and for females. Reports are known of alligators in China reaching in past centuries, but these are no longer thought to be accurate. The largest reported female measured and weighed , while the largest reported male measured and weighed .
The Chinese alligator is almost completely black or dark gray in color as an adult. It has a short and broad snout, which points slightly upwards and narrows at the end. Its head is robust, more so than that of the American alligator, with a bony septum dividing its nostrils. It has 72–76 teeth, of which 13–14 are maxillary, five premaxillary, and 18–19 mandibular. Four specimens measuring in length and weighing had a bite force of . Unlike the American alligator, the Chinese alligator is fully armored, including its belly. It contains up to 17 rows of scales across its body, which are soft on its belly and side and rougher on its back. Its upper eyelids have bony plates on them, a feature usually not present in the American alligator. Its tail is wider than that of the American alligator. It does not have webbed feet, in contrast to the American alligator, which has extensive webbing on its toes.
Ecology
The Chinese alligator brumates in burrows during winter. After this period of dormancy, it frequently spends time in the sun before summer begins. It is nocturnal throughout summer, feeding at night and sheltering in the daytime, to avoid both humans and the summer heat. This behavior gives it the ability to live in areas where humans are common. A docile species, it generally does not intentionally hurt humans.
Burrowing
This alligator brumates from late October to mid-April, emerging in early May. It constructs its burrows next to ponds and other small bodies of water, using its head and front legs to dig into the ground. They can be large and complex, containing multiple rooms, water pools, and entrances. Most of them are long, with each room having enough space for alligators to turn around after entering. Outside of winter, the burrows serve as retreat sites for the alligators and in summer are where they take shelter in the daytime. The temperature inside them is never colder than . The burrows can be problematic for farmers, as they cause destruction of farm dykes.
Life cycle
The breeding season of the Chinese alligator is early summer, with the rate of mating being highest in mid-June. The alligator breeds earlier in the year if temperatures are higher. During the time of mating, males commonly search around ponds to find a mate and both male and female specimens are often aggressive to each other. The species exhibits polygamy, with single males mating with multiple females and/or a single female mating with several males. A study of 50 clutches showed multiple paternity in 60% of them, with up to three males contributing. Nests are typically built about 2–3 weeks after mating, from July to late August. Constructed by the females, they are composed of rotting plants, such as leaves, and are high. Females prefer to assemble them in areas that have a thick canopy and are far from human disturbance. Because islands frequently satisfy both of these conditions, they are often used as nesting sites. Nests are always near water sources. Individuals often return to the same nesting site yearly, although intraspecific competition and environmental changes can force them to change nesting sites.
Generally laid at night, mating typically produces 20–30 eggs, although according to the International Union for Conservation of Nature (IUCN), clutch size ranges between 10 and 40 eggs. After the eggs are laid, the females sometimes leave the nest, but other times stay to protect the eggs. The eggs are about in length, in diameter, and in weight, making them smaller than the eggs of any other crocodilian. They are typically incubated for about 70 days. On average, the temperature of incubation is , including the day and night. This temperature controls whether a young alligator will be male or female (temperature-dependent sex determination), a feature present in many other reptiles. A higher incubation temperature also increases the hatching rate. Young hatch in September, assisted by their mothers.
Newborn alligators, like their eggs, are the smallest of any crocodilian, with a length of and weight of . Unlike adults, they have light speckles on their bodies and heads. Mothers help them leave the nest and bring them to the water after hatching. They grow very little in their first year, due to being able to feed for only about 2 months after hatching before the winter. A 2002 study showed that the Chinese alligator is two-thirds the length of the American alligator and one-half its weight at birth, but is one-half its length and one-tenth its weight after one year. Young depend on their mothers to protect them during their first winter, as their small size makes them an easy prey target.
The alligator grows quickly in its first few years, with its growth rate slowing at age five. According to the National Zoological Park, females reach maturity roughly four to five years after birth, although other sources estimate that they mature at age six to seven. It can live to over 50 years, and has been known to reach age 70 in captivity. It cannot breed past its 50s.
Feeding
The Chinese alligator is an opportunistic feeder, meaning that it can prey on a variety of different animals depending on what is available. It is a carnivore, mostly eating fish and invertebrates, such as crustaceans, insects, mussels, clams, and snails. When possible, it eats rodents, other small mammals, and aquatic birds as well. It has dull teeth, which allow it to eat prey with shells more easily. There is some speculation that they may prey on turtles as well. A study of the alligator in 1985 showed that snails were the most common animal in its diet at 63%, with 65% of that being river snails and 35% spiral-shelled snails. According to the survey, its diet also contained 16% rabbits, 8.3% mollusks, and 4.1% shrimp, with the remaining 6.8% being frogs, fish, and insects.
Vocalization
The Chinese alligator is a vocal species, making many different sounds in multiple situations. When communicating with nearby alligators, it produces sounds such as head slapping, hissing, and whining, which have a low sound pressure level (SPL). To communicate long-distance, it produces bellows, which have a high SPL. All of these sounds have a low frequency of less than 500 hertz, due to the alligator's densely vegetated habitat, which allow the sounds to spread across a greater area.
Both sexes participate in bellowing choruses during the mating season as adults. Lasting an average of 10 minutes, the alligators remain still for the entirety of the chorus, with both sexes responding equally in rough unison. The main purpose of these bellows is to call out to alligator specimens to collect at a specific pond, where individuals choose mates and engage in copulation. Alligators may also bellow to publicize their size, a behavior which occurs in multiple other vertebrates. The size of a specimen is a significant factor for mating; females only mate with males larger than themselves. Bellowing is most common at 6:00–7:00 am and 11:00–12:00 am CST. Although these bellows occur most frequently during the mating season, adults also bellow throughout the rest of the year.
Young Chinese alligators often communicate with each other and their parents using vocal signals to "maintain group cohesion". Young also make sounds when in danger, which alert adults to help and caution nearby young of the threat. Embryos produce distinctive sounds inside their eggs, which alert the adult female that the nest is ready to be opened. These vocalizations are high-pitched, while their danger calls are louder.
Distribution and population
The modern range of the Chinese alligator is extremely restricted. Historically, the alligator was widely distributed in the Yangtze river system. By the late 1980s, it was restricted to small ponds in six counties in the southeastern part of Anhui province; and since of 2015, the IUCN estimates the area of occupation at about . It is the only species in the family Alligatoridae that lives in the Eastern Hemisphere.
Habitat
The habitat of the Chinese alligator is bodies of fresh water, particularly wetlands and ponds, in areas transitioning between subtropical and temperate climates. It lives at the base of mountains, in areas where grass and shrubs are common. Habitat loss has also forced it to live at higher elevations than it prefers, where the weather is colder and the soil is unfit for burrow digging. Crocodilian conservationist John Thorbjarnarson observed a female who had to build her nest of pine needles rather than the usual plants; the eggs died due to the pine needles not being able to warm them properly.
Population and range trend
The oldest record of the Chinese alligator is a skeleton fragment found in western Japan. The fossil is estimated to be from the late Pliocene period, 3 million years ago (Mya). The skeleton showed that the species was larger at the time than it is currently, with a total length of at least . Alligators are believed to have moved into various parts of Japan either before 25 Mya or after 10 Mya and were extirpated from there during the Plio-Pleistocene period, due to Japan's increased isolation from the continent and harsh climate conditions.
The population of the Chinese alligator began to decline in 5000 BC, when human civilization started to grow in China, after having been very abundant in the lower Yangtze area. This area was one of the first places in the world to farm rice, causing much of the alligator's habitat to be destroyed in favor of rice farms. In the 1700s, much of the Chinese alligator's habitat was replaced with farming fields after a large number of people had moved into the area. By the 20th century, its range was reduced to a few small areas around the Yangtze. In the 1950s, the alligator was in three distinct areas: the southern area of the Yangtze (Chang Jiang) from Pengze to the western shore of Lake Tai (Tai Hu), the mountainous regions of southern Anhui, and the provinces of Jiangsu and Zhejiang, primarily in lakes, streams, and marshes. By the 1970s, it was restricted to small parts of southern Anhui and Zhejiang, at which time the population was about 1,000.
In 1998, the population of the Chinese alligator was the lowest it had ever been; the largest area it lived in was a small pond along the Yangtze surrounded by farmland, which held 11 alligators. In 1999, the Wildlife Conservation Society estimated that 130–150 individuals were left in the wild. According to The New York Times, the population was less than 130 in 2001; at this time, alligators sometimes wandered around to look for a suitable habitat, but were unsuccessful due to their habitat having been turned into rice fields. In 2003, the population began to gradually increase after having been roughly stable between 1998 and 2003. A survey of the population by the Anhui National Nature Reserve for Chinese Alligator (ANNRCA) in 2005 deduced that between 92 and 114 adults and 66 young remained in the wild. The survey reasoned that the species' population was growing in four sites, but stable in the rest of the alligator's range. A 2012 journal article estimated the population at the time to be 120–150. A 2015 survey observed 64 individuals, of which 32 were adults, estimating that the total number of adults was 68–86 and the total population 136–173. Wang Renping, the head of the ANNRCA, stated in 2017 that about 300 specimens existed in the wild, some of which had been born captive and reintroduced to the wild. As of 2018 the population is not considered to be further declining.
Due to the low wild population of the Chinese alligator, high inbreeding is a major concern threatening their chances for long-term survival.
Reasons for population decline
Considered to be one of the most endangered crocodilians in the world, the Chinese alligator's biggest threats in the late 20th century were human killing and habitat loss. A majority of the species' wetland habitats were destroyed to construct rice paddies and dams. During the 1970s and 1980s, humans sometimes killed the alligators, because they believed they were pests, out of fear, or for their meat. Their meat was thought to have the ability to cure colds and prevent cancer and their organs were sold for medicinal purposes. In several restaurants and food centers in China's more prosperous areas, young alligators were allowed to roam free with their mouths taped shut, and were subsequently killed for human consumption, served as a dish of rice, vegetables, and chopped up alligator flesh. In the late 20th century, people living in the range of the Chinese alligator ate its meat due to believing that it was dragon meat.
The Yangtze was flooded in the winter of 1957, which is believed to have caused many Chinese alligators to drown. Rats, which this species eat, have been poisoned by farmers, so were also a cause for the diminishing of the species. The organochlorine compound sodium pentachlorophenate was used to kill snails in agricultural fields starting in 1958, which incidentally poisoned the alligators as well. Other factors that led to the endangerment of the alligator include natural disasters and geographic separation.
Status and conservation
In its native country, the Chinese alligator has been listed as a Class I endangered species since 1972, which gives it the highest possible degree of legal protection and makes killing or capturing the species in the wild forbidden. It is listed as a CITES Appendix I species and an endangered species by the U.S. Fish & Wildlife Service. Following six assessments as endangered from 1982 to 1994, it is classified as critically endangered on the IUCN Red List as of 2017. In 1982, the Anhui National Nature Reserve for Chinese Alligator (ANNRCA) was created, a reserve spanning across the entire range of the Chinese alligator, now covering an area of .
In captivity
As of 2016, at least 20,000 Chinese alligators are living in captivity due to captive-breeding programs, the first initiated in the 1970s. Captive-born Chinese alligators have been reintroduced into their native range, boosting the wild population. Six specimens were released from captivity in 2007, followed by six more in June 2015. As of June 2016, the largest group of Chinese alligators to have been released in the wild was when 18 specimens were reintroduced to Langxi County, part of the species' native habitat, on May 22, 2016. These releases have proven successful, with individuals adapting well to a life in the wild and breeding. A year after the 2007 release, 16 young alligators were found living in the wild. 60 alligator eggs were observed in 2016, distributed in three nests at a wetland park. Although a typhoon in September the same year flooded and eliminated two of the nests, three hatchlings were found in the same area several days after.
China
The two largest breeding centers for the Chinese alligator are in, or near, the areas where Chinese alligators are still found in the wild. The Anhui Research Center for Chinese Alligator Reproduction (ARCCAR) is the largest of them, housing roughly 15,000 Chinese alligators as of 2016. The center is from the city of Xuancheng, where it makes use of a series of ponds in a small valley. Founded in 1979, the ARCCAR was stocked with 212 alligators collected from the wild over the first decade after its establishment, and received alligator eggs collected by the area's residents and the ARCCAR's own staff from the nests of wild alligators as well. In 1988, the first eggs by human-bred alligators were laid. The reserve decided to reintroduce some of its alligators in the wild in 2001, which was carried out in 2003 when three alligators were released. The alligator breeding was so successful that the ARCCAR began to use the alligators for local meat consumption and live animals for the European pet market, with the profits from these activities continuing to fund the breeding centers.
The other major breeding center for the species is the Changxing Chinese Alligator Nature Reserve (CCANR) or Changxing Nature Reserve and Breeding Center for Chinese Alligators (CNRBRCCA), in Changxing County, Zhejiang, about east of the ARCCAR. Originally known as the Yinjiabian Alligator Conservation Area (尹家边扬子鳄保护区), the breeding center was established in 1982. Unlike the ARCCAR, where alligator eggs are collected by the center's staff for incubation in controlled condition, the CCANR allows eggs to hatch naturally. According to a 2013 official report, the CCANR housed almost 4,000 alligators, including 2,089 young (1–3 years old), 1,598 juveniles (4–12 years old), and 248 adults (13+ years old). By 2016, 5,500 specimens were housed at the center.
In 2003, the ARCCAR received a donation of $1.2 million from the State Forestry and Grassland Administration of China (SFGA) and $740,000 from the government of Anhui. This allowed the organization to create two new breeding areas to hold the alligators, each, as well as heighten the existing fence. The same year, the CCANR received a donation of $600,000 from the SFGA and $800,000 from the government of Changxing, enabling it to reinstate wetlands for the alligators and enhance its facilities. Both the ARCCAR and the CCANR position themselves as tourist attractions, where paying visitors can view alligators and learn about them.
Multiple other breeding facilities that house the Chinese alligator exist in various provinces of China, as well as private breeding farms and museums.
Foreign countries
The Chinese alligator is also kept and bred at many zoos and aquariums in North America and Europe. Some individuals bred there have been returned to China for reintroduction to the wild. The first time the alligators were ever transported internationally is believed to have been when several were taken from China to the United States in the 1950s. In November 2017, four Chinese alligators were transported from their natural habitat in China to Shizuoka, Japan.
Among the North American zoos and aquariums keeping this species are the Bronx Zoo, Cincinnati Zoo, Great Plains Zoo, Sedgwick County Zoo, Philadelphia Zoo, San Diego Zoo, Santa Barbara Zoo, Smithsonian National Zoological Park, and St. Louis Zoo. In Europe, about 25 zoos and aquariums keep the species, such as the Barcelona Zoo (Spain), Parque de las Ciencias (Granada) (Spain), Bioparco di Roma (Italy), Crocodile Zoo (Denmark), Moscow Zoo (Russia), Pairi Daiza (Belgium), Paradise Wildlife Park (England), Parken Zoo (Sweden), Prague Zoo (Czech Republic), Tallinn Zoo (Estonia) and Tierpark Berlin (Germany).
Chinese dragon association
Some writers have suggested that the Chinese alligator was the inspiration for the Chinese dragon. This theory was widespread in the early 1900s, and the idea was later revisited by John Thorbjarnarson and Xiaoming Wang. According to The New York Times, the association with the "beneficent" mythological creature is an advantage for the species.
Unlike dragons in myths of the Western Hemisphere, the Chinese dragon is portrayed as a symbol of "royal power and good fortune", frequently helping and saving people. It is able to swim in water or air. The relatively harmless nature of the Chinese alligator is believed to have been an influence for the helpful nature of the dragon. The fact that the alligator ends its brumation when the rainy season begins and returns to its burrows when the rainwater in rivers recedes, as well as the fact that it lives in bodies of water, may be the reason for the dragon's portrayal as a water-related mythological creature. Alligator drums may have been used to simulate the species' vocalizations during the mating season, which humans associated with the dragon's "power of summoning rainclouds".
| Biology and health sciences | Crocodilia | Animals |
1246718 | https://en.wikipedia.org/wiki/Organic%20matter | Organic matter | Organic matter, organic material, or natural organic matter refers to the large source of carbon-based compounds found within natural and engineered, terrestrial, and aquatic environments. It is matter composed of organic compounds that have come from the feces and remains of organisms such as plants and animals. Organic molecules can also be made by chemical reactions that do not involve life. Basic structures are created from cellulose, tannin, cutin, and lignin, along with other various proteins, lipids, and carbohydrates. Organic matter is very important in the movement of nutrients in the environment and plays a role in water retention on the surface of the planet.
Formation
Living organisms are composed of organic compounds. In life, they secrete or excrete organic material into their environment, shed body parts such as leaves and roots and after organisms die, their bodies are broken down by bacterial and fungal action. Larger molecules of organic matter can be formed from the polymerization of different parts of already broken down matter. The composition of natural organic matter depends on its origin, transformation mode, age, and existing environment, thus its bio-physicochemical functions vary with different environments.
Natural ecosystem functions
Organic matter is common throughout the ecosystem and is cycled through decomposition processes by soil microbial communities that are crucial for nutrient availability. After degrading and reacting, it can move into soil and mainstream water via waterflow. Organic matter provides nutrition to living organisms. Organic matter acts as a buffer in aqueous solutions to maintain a neutral pH in the environment. The buffer acting component has been proposed to be relevant for neutralizing acid rain.
Source cycle
Some organic matter not already in the soil comes from groundwater. When the groundwater saturates the soil or sediment around it, organic matter can freely move between the phases. Groundwater has its own sources of natural organic matter including:
organic matter deposits, such as kerogen and coal.
soil and sediment organic matter.
organic matter infiltrating into the subsurface from rivers, lakes, and marine systems."
Organisms decompose into organic matter, which is then transported and recycled. Not all biomass migrates, some is rather stationary, turning only over the course of millions of years.
Soil organic matter
The organic matter in soil derives from plants, animals and microorganisms. In a forest, for example, leaf litter and woody materials fall to the forest floor. This is sometimes referred to as organic material. When it decays to the point in which it is no longer recognizable, it is called soil organic matter. When the organic matter has broken down into a stable substance that resists further decomposition it is called humus. Thus soil organic matter comprises all of the organic matter in the soil exclusive of the material that has not decayed.
An important property of soil organic matter is that it improves the capacity of a soil to hold water and nutrients, and allows their slow release, thereby improving the conditions for plant growth. Another advantage of humus is that it helps the soil to stick together which allows nematodes, or microscopic bacteria, to easily decay the nutrients in the soil.
There are several ways to quickly increase the amount of humus. Combining compost, plant or animal materials/waste, or green manure with soil will increase the amount of humus in the soil.
Compost: decomposed organic material.
Plant and animal material and waste: dead plants or plant waste such as leaves or bush and tree trimmings, or animal manure.
Green manure: plants or plant material that is grown for the sole purpose of being incorporated with soil.
These three materials supply nematodes and bacteria with nutrients for them to thrive and produce more humus, which will give plants enough nutrients to survive and grow.
Soil organic matter is crucial to all ecology and to all agriculture, but it is especially emphasized in organic farming, where it is relied upon especially heavily.
Priming effect
The priming effect is characterized by intense changes in the natural process of soil organic matter (SOM) turnover, resulting from relatively moderate intervention with the soil. The phenomenon is generally caused by either pulsed or continuous changes to inputs of fresh organic matter (FOM). Priming effects usually result in an acceleration of mineralization due to a trigger such as the FOM inputs. The cause of this increase in decomposition has often been attributed to an increase in microbial activity resulting from higher energy and nutrient availability released from the FOM. After the input of FOM, specialized microorganisms are believed to grow quickly and only decompose this newly added organic matter. The turnover rate of SOM in these areas is at least one order of magnitude higher than the bulk soil.
Other soil treatments, besides organic matter inputs, which lead to this short-term change in turnover rates, include "input of mineral fertilizer, exudation of organic substances by roots, mere mechanical treatment of soil or its drying and rewetting."
Priming effects can be either positive or negative depending on the reaction of the soil with the added substance. A positive priming effect results in the acceleration of mineralization while a negative priming effect results in immobilization, leading to N unavailability. Although most changes have been documented in C and N pools, the priming effect can also be found in phosphorus and sulfur, as well as other nutrients.
Löhnis was the first to discover the priming effect phenomenon in 1926 through his studies of green manure decomposition and its effects on legume plants in soil. He noticed that when adding fresh organic residues to the soil, it resulted in intensified mineralization by the humus N. It was not until 1953, though, that the term priming effect was given by Bingeman in his paper titled, The effect of the addition of organic material on the decomposition of an organic soil. Several other terms had been used before priming effect was coined, including priming action, added nitrogen interaction (ANI), extra N and additional N. Despite these early contributions, the concept of the priming effect was widely disregarded until about the 1980s-1990s.
The priming effect has been found in many different studies and is regarded as a common occurrence, appearing in most plant soil systems. However, the mechanisms which lead to the priming effect are more complex than originally thought, and still remain generally misunderstood.
Although there is a lot of uncertainty surrounding the reason for the priming effect, a few undisputed facts have emerged from the collection of recent research:
The priming effect can arise either instantaneously or very shortly (potentially days or weeks) after the addition of a substance is made to the soil.
The priming effect is larger in soils that are rich in C and N as compared to those poor in these nutrients.
Real priming effects have not been observed in sterile environments.
The size of the priming effect increases as the amount of added treatment to the soil increases.
Recent findings suggest that the same priming effect mechanisms acting in soil systems may also be present in aquatic environments, which suggests a need for broader considerations of this phenomenon in the future.
Decomposition
One suitable definition of organic matter is biological material in the process of decaying or decomposing, such as humus. A closer look at the biological material in the process of decaying reveals so-called organic compounds (biological molecules) in the process of breaking up (disintegrating).
The main processes by which soil molecules disintegrate are by bacterial or fungal enzymatic catalysis. If bacteria or fungi were not present on Earth, the process of decomposition would have proceeded much slower.
Various factors impact the decomposition of organic matter including its chemical properties and other environmental parameters. Metabolic capabilities of the microbial communities play a crucial role on decomposition since they are highly connected with the energy availability and processing. In terrestrial ecosystems the energy status of soil organic matter has been shown to affect microbial substrate preferences. Some organic matter pools may be energetically favorable for the microbial communities resulting in their fast oxidation and decomposition, in comparison with other pools where microbial degraders get less return from the energy they invest. By extension, soil microorganisms preferentially mineralize high-energy organic matter, avoiding decomposing less energetically dense organic matter.
Organic chemistry
Measurements of organic matter generally measure only organic compounds or carbon, and so are only an approximation of the level of once living or decomposed matter. Some definitions of organic matter likewise only consider "organic matter" to refer to only the carbon content or organic compounds and do not consider the origins or decomposition of the matter. In this sense, not all organic compounds are created by living organisms, and living organisms do not only leave behind organic material. A clam's shell, for example, while biotic, does not contain much organic carbon, so it may not be considered organic matter in this sense. Conversely, urea is one of many organic compounds that can be synthesized without any biological activity.
Organic matter is heterogeneous and very complex. Generally, organic matter, in terms of weight, is:
45–55% carbon
35–45% oxygen
3–5% hydrogen
1–4% nitrogen
The molecular weights of these compounds can vary drastically, depending on if they repolymerize or not, from 200 to 20,000 amu. Up to one-third of the carbon present is in aromatic compounds in which the carbon atoms form usually six-membered rings. These rings are very stable due to resonance stabilization, so they are challenging to break down. The aromatic rings are also susceptible to electrophilic and nucleophilic attacks from other electron-donating or electron-accepting material, which explains the possible polymerization to create larger molecules of organic matter.
Some reactions occur with organic matter and other materials in the soil to create compounds never seen before. Unfortunately, it is challenging to characterize these because so little is known about natural organic matter in the first place. Research is currently being done to determine more about these new compounds and how many are being formed.
Aquatic
Aquatic organic matter can be further divided into two components: (1) dissolved organic matter (DOM), measured as colored dissolved organic matter (CDOM) or dissolved organic carbon (DOC), and (2) particulate organic matter (POM). They are typically differentiated by that which can pass through a 0.45 micrometre filter (DOM), and that which cannot (POM).
Detection
Organic matter is important in water and wastewater treatment and recycling, natural aquatic ecosystems, aquaculture, and environmental rehabilitation. It is, therefore, important to have reliable methods of detection and characterisation, for both short- and long-term monitoring. Various analytical detection methods for organic matter have existed for up to decades to describe and characterise organic matter. These include, but are not limited to: total and dissolved organic carbon, mass spectrometry, nuclear magnetic resonance (NMR) spectroscopy, infrared (IR) spectroscopy, UV-Visible spectroscopy, and fluorescence spectroscopy. Each of these methods has its advantages and limitations.
Water purification
The same capability of natural organic matter that helps with water retention in the soil creates problems for current water purification methods. In water, organic matter can still bind to metal ions and minerals. The purification process does not necessarily stop these bound molecules but does not cause harm to any humans, animals, or plants. However, because of the high reactivity of organic matter, by-products that do not contain nutrients can be made. These by-products can induce biofouling, which essentially clogs water filtration systems in water purification facilities, as the by-products are larger than membrane pore sizes. This clogging problem can be treated by chlorine disinfection (chlorination), which can break down residual material that clogs systems. However, chlorination can form disinfection by-products.
Water with organic matter can be disinfected with ozone-initiated radical reactions. The ozone (three oxygens) has powerful oxidation characteristics. It can form hydroxyl radicals (OH) when it decomposes, which will react with the organic matter to shut down the problem of biofouling.
Vitalism
The equation of "organic" with living organisms comes from the now-abandoned idea of vitalism, which attributed a special force to life that alone could create organic substances. This idea was first questioned after Friedrich Wöhler artificially synthesized urea in 1828.
| Biology and health sciences | Biology basics | Biology |
1246810 | https://en.wikipedia.org/wiki/Humidifier | Humidifier | A humidifier is a household appliance or device designed to increase the moisture level in the air within a room or an enclosed space. It achieves this by emitting water droplets or steam into the surrounding air, thereby raising the humidity.
In the home, point-of-use humidifiers are commonly used to humidify a single room, while whole-house or furnace humidifiers, which connect to a home's HVAC system, provide humidity to the entire house. Medical ventilators often include humidifiers for increased patient comfort. Large humidifiers are used in commercial, institutional, or industrial contexts, often as part of a larger HVAC system.
Overview
Humidification calculation
Humidity per hour: X = Air changes per hour (ACPH) * M³ * density of air * humidity ratio
Humidity per day: X * 24
The air changes per hour (ACPH) ranges wildly based on:
Ventilation: Values may be obtained from the HVAC maintainer that routinely (typically every third year or so) tests the ventilation of the residence.
Insulation leakage: Measured with a standard blower door test.
Cubic meters: The volume of the room, excluding the bathroom that should be kept closed since it ventilates humidity.
Density of air: Typically 1.2 for dry air.
Humidity
Current relative humidity: 20%
Humidity needed to reach 55%: 35%
Humidity ratio for 35%: 0.0051
For example, a typical modern apartment of with closed windows (wood isolation) may consume to raise the relative humidity from 20% to 55%: 24 * 1,53 L/h (2*1×125×1.2×0.0051):
Air changes per hour
Ventilation: 2
Insulation leakage: 1 (a few windows)
Cubic meters: 125 m3 (50 m2 * 2.5 m height)
Density of air: Typically 1.2 for dry air.
Humidity
Current relative humidity: 20%
Humidity needed to reach 55%: 35%
Humidity ratio for 35%: 0.0051
Prevention of low indoor humidity
Low humidity may occur in hot, dry desert climates, or indoors in artificially heated spaces. In winter, especially when cold outside air is heated indoors, the humidity may drop to as low as 10–20%. A relative humidity of 30% to 50% is recommended for most homes.
Health treatment
Prevention of dermatitis: Low humidity can cause adverse health effects and may cause atopic dermatitis, and seborrhoeic dermatitis.
Management of hair loss: Commonly, patients with seborrhoeic dermatitis experience mild redness, scaly skin lesions and in some cases hair loss.
Prevention of dry mucous membranes and cough: By drying out mucous membranes such as the lining of the nose and throat, it may also lead to a snoring problem, and can cause respiratory distress.
Prevention of dry eye syndrome.
Improved apparent temperature: The heat index and humidex measure the effect of humidity on the perception of temperatures above . In humid conditions, the air feels much hotter, because less perspiration evaporates from the skin.
Improved climate for material
Low humidity can affect wooden furniture, causing shrinkage and loose joints or cracking of pieces. Books, papers, and artworks may shrink or warp and become brittle in very low humidity.
In addition, static electricity may become a problem in conditions of low humidity, destroying semiconductor devices, causing static cling of textiles, and causing dust and small particles to stick stubbornly to electrically charged surfaces.
Negative impact of overuse of humidifiers
An indoor relative humidity of less than 51% resulted in significant reductions in mite and allergen levels. Overuse of a humidifier can raise the relative humidity to excessive levels, promoting the growth of dust mites and mold, and can also cause hypersensitivity pneumonitis (humidifier lung). A properly installed and located hygrostat should be used to monitor and control humidity levels automatically, or a well-informed and conscientious human operator must constantly check for correct humidity levels.
A humidity below 50% can prevent water condensation on building materials.
A dehumidifier can be used to balance the humidity.
Humidifiers
Industrial humidifiers are used when a specific humidity level must be maintained to prevent static electricity buildup, preserve material properties, and ensure a comfortable and healthy environment for workers or residents.
Static problems are prevalent in industries such as packaging, printing, paper, plastics, textiles, electronics, automotive manufacturing and pharmaceuticals. Friction can produce static buildup and sparks when humidity is below 45% relative humidity (RH). Between 45% and 55% RH, static builds up at reduced levels, while humidity above 55% RH ensures that static will never build up. The American Society of Heating, Refrigerating and Air Conditioning Engineers (ASHRAE) has traditionally recommended a range of 45–55% RH in data centers to prevent sparks that can damage IT equipment. Humidifiers are also used by manufacturers of semiconductors and in hospital operating rooms.
Printers and paper manufacturers use humidifiers to prevent shrinkage and paper curl. Humidifiers are needed in cold storage rooms to preserve the freshness of food against the dryness caused by hot temperatures. Art museums use humidifiers to protect sensitive works of art, especially in exhibition galleries, where they combat the dryness caused by heating for the comfort of visitors during winter.
Natural humidifiers
Natural humidifiers don't use or need a demineralization filter because the water is slowly evaporated which leaves the mineral deposit at the bottom of the container. However, natural humidifiers raise the humidity very slowly even if their water surface area is large.
Common sources
Human water losses, both respiratory, and insensible water loss like sweat, range in average 0.75 L/d in sedentary adults. However, most people do not spend most of the day at home.
Houseplants may also be used as natural humidifiers, especially if they are placed in fabric flowerpots, since they evaporate water into the air through transpiration. Care must still be taken to prevent bacteria or mold in the soil from growing to excessive levels, or from dispersing into the air. The presence of sciarids (like fungus gnats) in houseplants may indicate overwatering.
Hanging laundry will increase the humidity.
Homemade
One type of evaporative humidifier makes use of just a reservoir and wick. Sometimes called a "natural humidifier", these are usually non-commercial devices that can be assembled at little or no cost.
One version of a natural humidifier uses a stainless steel bowl, partially filled with water, covered by a towel. A waterproof weight is used to sink the towel in the center of the bowl. There is no need for a fan, because the water spreads through the towel by capillary action and the towel surface area is large enough to provide for rapid evaporation. The stainless steel bowl is much easier to clean than typical humidifier water tanks. This, in combination with daily or every other day replacement of the towel and periodic laundering, can control the problem of mold and bacteria.
Electric humidifiers
Evaporative humidifiers
An "evaporative", "cool moisture", or "wick humidifier", consists of just three basic parts: a reservoir, a wick, and a fan.
The wick is made of a porous material that absorbs water from the reservoir and provides a larger surface area for it to evaporate from. The fan is adjacent to the wick and blows air onto the wet wick to aid in the evaporation of the water. Evaporation from the wick is dependent on relative humidity. A room with low humidity will have a higher evaporation rate compared to a room with high humidity. Therefore, this type of humidifier is partially self-regulating; as the humidity of the room increases, the water vapor output naturally decreases.
These wicks become moldy if they are not dried out completely between fillings, and become saturated with mineral deposits over time. They regularly need rinsing or replacement; if this does not happen, air cannot pass through them, and humidifier stops humidifying the area it is in and the water in the tank remains at the same level.
Evaporative humidifiers function similarly to evaporative coolers.
Impeller humidifiers
Impeller humidifiers (a type of cool mist humidifier) are usually noisier than others. It uses a rotating disc to fling water at a diffuser, which breaks the water into fine droplets that float into the air. The water supply must be kept scrupulously clean, or there is a risk of spreading bacteria or mold into the air.
Ultrasonic humidifiers
An ultrasonic humidifier uses a ceramic diaphragm vibrating at an ultrasonic frequency to create water droplets that silently exit the humidifier in the form of cool fog. Usually the mist gets forced out by a tiny fan, while some ultra mini models have no fans. The models without fans are meant mainly for personal use. Ultrasonic humidifiers use a piezoelectric transducer to create a high frequency (1-2 MHz) mechanical oscillation in a film of water. This forms an extremely fine mist of droplets about one micron in diameter, that is quickly evaporated into the air flow.
Unlike the humidifiers that boil water, these water droplets will contain any impurities that are in the reservoir, including minerals from hard water (which then forms a difficult-to-remove sticky white dust on nearby objects and furniture). Any pathogens growing in the stagnant tank will also be dispersed in the air. Ultrasonic humidifiers should be cleaned regularly to prevent bacterial contamination from being spread throughout the air.
The amount of minerals and other materials can be greatly reduced by using distilled water. Special disposable demineralization cartridges may also reduce the amount of airborne material, but the EPA warns, "the ability of these devices to remove minerals may vary widely." The mineral dust may have negative health effects. Wick humidifiers trap the mineral deposits in the wick; vaporizer types tend to collect minerals on or around the heating element and require regular cleaning with vinegar or citric acid to control buildup.
Steam Humidifiers
Steam humidifiers, or warm mist humidifiers, are equipped with a heating element.
A medicated inhalant can also be added to the steam vapor to help reduce cough. Vaporizers may be healthier than cool mist types of humidifiers because steam is less likely to convey mineral impurities or microorganisms from the standing water in the reservoir. However, boiling water requires significantly more energy than other techniques. The heat source in poorly designed humidifiers can overheat, causing the product to melt, leak, and start fires.
Tanks
The water is usually supplied by manually filling the unit on a periodic basis.
Top fill: A top fill tank has a hole at the top to make it convenient to refill the tank on daily basis.
Bottom fill: Bottom filled humidifiers has removable water tank, often attached with a replaceable demineralization filter. Both the cap and the filter are often universal, which means that they can be switched between many tank cap humidifiers.
Tank cap-based humidifiers are more inconvenient to fill than top fill tanks, because the demineralization filter has to be unscrewed and screwed back on a daily basis. Also, water drops from the tank attachment may get in contact with the hygrometer, which will make it harder to operate properly if they are not operated with caution.
Cleaning
Ultrasonic wave nebulizers trap the mineral deposits over time, even if filters are used, and require cleaning with vinegar or citric acid to control buildup. It is easy to remove the buildup from the nebulizer in humidifiers with removable base parts. Other models must be screwed. Some models require monthly maintenance, while other models may run for years unmaintained.
Activated charcoal filter
Some humidifiers have an activated charcoal filter to reduce contaminants.
Demineralization filter
Humidity with mineral deposits build up on furniture, and static objects like computer monitors, TVs, etc.
Most models with heating elements come with replacement demineralization filters (cartridges, or plastic foam, for the tank, and fabric near the heater in the reservoir) to reduce the amount of minerals entering the system. The cost of these filters is about $10–15 per cartridge or 10-pack for plastic foam, or fabric filters. Cartridges may last for 6 months, and plastic foam/fabric filters may last one month, if the humidifier is frequently used. Cartridges can be used for most humidifiers, but some models have a proprietary format.
Some humidifiers (for example, some cold mist humidifiers) don't come with any filters at all.
Clogged filters should be replaced when they cannot deliver water, which is indicated with a "no water" icon/alert signal in modern humidifiers. A manual way to figure out if a filter is clogged is to shake it, if the content of it (even if wet) doesn't move around, then it's unlikely that sufficient water will pass through.
Ultraviolet germicidal irradiation (UVGI)
Some humidifiers have an integrated ultraviolet germicidal irradiation (UVGI) feature.
Water ionizer
Some humidifiers have an integrated water ionizer.
Maintenance
A torx screwdriver is often needed to open consumer electronics.
Humidifiers with a centrifugal fan may become noisy due to build up which cannot effectively be removed with vinegar for example. However, the fan can be replaced, and the model name of the fan and its specifications can be figured out by disassembling the humidifier and detaching it.
Fixed-installation humidifiers
For buildings with a forced-air furnace, a humidifier may be installed into the furnace. They can also protect wooden objects, antiques and other furnishings which may be sensitive to damage from overly dry air. In colder months, they may provide modest energy savings, since as humidity increases, occupants may feel warm at a lower temperature.
Bypass humidifiers are connected between the heated and cold air return ducts, using the pressure difference between these ducts to cause some heated air to make a bypass through the humidifier and return to the furnace.
Any humidifiers should usually be disabled during the summer months if air conditioning is used; air conditioners partially function to reducing indoor humidity, and having a humidifier continue to operate will waste significant amounts of energy.
Drums
Drum style (bypass) uses a pipe to bring water directly to a reservoir (a pan) attached to the furnace. The water level in the pan is controlled by a float valve, similar to a small toilet tank float. The wick is typically a foam pad mounted on a drum and attached to a small motor; hot air enters the drum at one end and is forced to leave through the sides of the drum. When the hygrostat calls for humidity, the motor is turned on causing the drum to rotate slowly through the pan of water and preventing the foam pad from drying out.
Advantages include:
Low cost
Inexpensive maintenance (drum-style pads are cheap and readily available)
Disadvantages include:
Requirement for frequent (approximately monthly) inspections of cleanliness and pad condition
Water evaporation even when humidification is not required (due to the pan of water which remains exposed to a high velocity air stream)
Mold growth in the pan full of water (this problem is exacerbated by the large quantity of air, inevitably carrying mold spores, passing through the humidifier whether in use or not).
For the latter reason especially, drum-style humidifiers should always be turned off at the water supply during summer (air conditioning) months, and should always be used with high quality furnace air filters (MERV ratings as high as possible to ensure small numbers of mold spores reaching the humidifier pan) when the water supply is turned on.
Disc wheels
A disc wheel style (bypass) is very similar in design to the drum style humidifiers; this type of furnace humidifier replaces the foam drumming with a number of plastic discs with small grooves on both sides. This allows for a very large evaporative surface area, without requiring a great deal of space. Unlike the drum style humidifiers, the disc wheel does not need regular replacement.
Advantages include:
Very low maintenance (basin of humidifier should be cleaned out periodically, unless an automatic flushing device is installed)
No regular replacement of parts necessary
Higher output due to large evaporative surface area
Can be installed in hard water situations
Maintains efficiency throughout lifespans
Disadvantages include:
Higher price
Water evaporation even when humidification is not required (due to the pan of water which remains exposed to a high velocity air stream)
Bypass flow-through
Bypass flow-through style (bypass – also known as "biscuit style" or many other, similar variant names) uses a pipe to bring water directly to an electrically controlled valve at the top of the humidifier. Air passes through an aluminum "biscuit" (often called a pad; the term "biscuit" emphasizes the solid rather than foamy form) which is similar to a piece of extremely coarse steel wool. The "biscuit" has a coating of a matte ceramic, resulting in an extremely large surface area within a small space. When the hygrostat calls for humidity, the valve is opened and causes a spray of water onto the "biscuit". Hot air is passed through the "biscuit", causing the water to evaporate from the pad and be carried into the building.
Advantages include:
Reduced maintenance (new "biscuit" is needed only when clogged with dust or mineral deposits, typically once per year)
Lack of a pan of potentially stagnant water to serve as a breeding ground for mold as with a drum-style humidifier
No incidental humidification caused by a constantly replenished pan of water in a high velocity air stream
Reduced requirement for expensive air filters
Uses little electricity
Disadvantages include:
A somewhat higher purchase price
Manufacturer and model-specific replacement biscuits (versus the relatively generic drum-style pads) may be more expensive and difficult to find
For most models, a portion of the water supplied to the unit is not evaporated. This can generate a considerable amount of waste water containing residual minerals, which does require connection to a drain. There is a limited selection of drain-less models that recirculate water, but mineral buildup must then be removed manually on a periodic basis.
Spray mist
Spray mist type uses a pipe, usually a small plastic one, to bring water directly to an electrically controlled valve (atomizer-this forces the water through a tiny orifice causing it to break up into tiny particles) in the humidifier. Water mist is sprayed directly into the supply air, and the mist is carried into the premises by the air flow.
Advantages include:
Simpler than bypass types to install, requiring a single cut hole for installation, no additional ducting.
Uses little electricity.
Small, compact unit which fits where other types cannot. (Approximately square.)
Because it does not require bypass ducting it does not undermine the pressure separation (and therefore, blower efficiency) of the return and supply ducts.
Does not require use of moisture pads (on-going expense).
Highly efficient usage of water. Does not generate waste water, and does not require separate connection to a drain.
Requires little maintenance. Periodic cleaning of nozzle may be required in hard water environments.
Lack of a pan of potentially stagnant water to serve as a breeding ground for mold as with a drum-style humidifier.
Disadvantages include:
Spray nozzle can become clogged in hard water situations, necessitating the use of water filter, periodic cleaning of nozzle, or nozzle replacement.
Disperses any minerals in the water into the airstream.
Invention
On July 27, 1962, Raymond Banks applied for a patent for a room humidifier. The patent was granted on November 3, 1964, to the inventor and originally assigned to Walton Labs, Inc.
Additional types
There are many different types of Plant and Home Humidifiers which can help maintain the humidity level of the home as required.
Additional types include non-bypass flow-through (fan augmented), steam, impeller or centrifugal atomizer, and under duct designs.
Problems
The U.S. EPA provides detailed information about health risks as well as recommended maintenance procedures. If the tap water contains a lot of minerals (also known as "hard water") then the ultrasonic or impeller humidifiers will produce a "white dust" (calcium is the most common mineral in tap water), which usually settles onto furniture, and is attracted to static electricity generating devices such as CRT monitors. The white dust can be prevented by using distilled water, or a demineralization cartridge in ultrasonic humidifiers. Bottled waters labeled "natural", "artesian", or "spring" may still have their original mineral content. The EPA reports that ultrasonic and impeller humidifiers spread the most mineral deposits and microorganisms, while evaporative and steam humidifers can allow the growth of microorganisms but generally disperse less of them into the air.,
In addition, a stuck or malfunctioning water supply valve can deliver large amounts of water, causing extensive water damage if undetected for any period of time. A water alarm, possibly with an automatic water shutoff, can help prevent this malfunction from causing major problems.
From 2006 to 2011, the disinfectant polyhexamethylene guanidine and other toxic materials were used as a cleaning agent for humidifier water tanks in Korea, leading to severe lung disease. Eighty children died of the disorder and nine adults either died or needed lung transplants. In the two years following the banning of cleaning chemicals for humidifier tanks, there were no new cases.
| Technology | Household appliances | null |
1247265 | https://en.wikipedia.org/wiki/Volume%20integral | Volume integral | In mathematics (particularly multivariable calculus), a volume integral (∭) is an integral over a 3-dimensional domain; that is, it is a special case of multiple integrals. Volume integrals are especially important in physics for many applications, for example, to calculate flux densities, or to calculate mass from a corresponding density function.
In coordinates
It can also mean a triple integral within a region of a function and is usually written as:
A volume integral in cylindrical coordinates is
and a volume integral in spherical coordinates (using the ISO convention for angles with as the azimuth and measured from the polar axis (see more on conventions)) has the form
Example
Integrating the equation over a unit cube yields the following result:
So the volume of the unit cube is 1 as expected. This is rather trivial however, and a volume integral is far more powerful. For instance if we have a scalar density function on the unit cube then the volume integral will give the total mass of the cube. For example for density function:
the total mass of the cube is:
| Mathematics | Multivariable and vector calculus | null |
1248929 | https://en.wikipedia.org/wiki/Fetal%20alcohol%20spectrum%20disorder | Fetal alcohol spectrum disorder | Fetal alcohol spectrum disorders (FASDs) are a group of conditions that can occur in a person who is exposed to alcohol during gestation. FASD affects 1 in 20 Americans, but is highly mis- and under-diagnosed.
The several forms of the condition (in order of most severe to least severe) are: fetal alcohol syndrome (FAS), partial fetal alcohol syndrome (pFAS), alcohol-related neurodevelopmental disorder (ARND), and neurobehavioral disorder associated with prenatal alcohol exposure (ND-PAE). Other terms used are fetal alcohol effects (FAE), partial fetal alcohol effects (PFAE), alcohol-related birth defects (ARBD), and static encephalopathy, but these terms have fallen out of favor and are no longer considered part of the spectrum.
Not all infants exposed to alcohol in utero will have detectable FASD or pregnancy complications. The risk of FASD increases with amount consumed, the frequency of consumption, and the longer duration of alcohol consumption during pregnancy, particularly binge drinking. The variance seen in outcomes of alcohol consumption during pregnancy is poorly understood. Diagnosis is based on an assessment of growth, facial features, central nervous system, and alcohol exposure by a multi-disciplinary team of professionals. The main criteria for diagnosis of FASD is nervous system damage and alcohol exposure, with FAS including congenital malformations of the lips and growth deficiency. FASD is often misdiagnosed as or comorbid with ADHD.
Almost all experts recommend that the mother abstain from alcohol use during pregnancy to prevent FASDs. As the woman may not become aware that she has conceived until several weeks into the pregnancy, it is also recommended to abstain while attempting to become pregnant. Although the condition has no known cure, treatment can improve outcomes. Treatment needs vary but include psychoactive medications, behavioral interventions, tailored accommodations, case management, and public resources.
Globally, one in 10 women drink alcohol during pregnancy, and the prevalence of having any FASD disorder is estimated to be at least 1 in 20. The rates of alcohol use, FAS, and FASD are likely to be underestimated, because of the difficulty in making the diagnosis and the reluctance of clinicians to label children and mothers. Some have argued that the FAS label stigmatizes alcohol use, while authorities point out that the risk is real. The condition has appeared in several works of fiction.
Signs and symptoms
The key signs of fetal alcohol syndrome (FAS) required for diagnosis include:
Growth deficiency or failure to thrive: slow fetal growth low birth weight or height, small head circumference (microcephaly)
Congenital malformations of lips: short palpebral fissure lengths, smooth philtrum, and thin upper lip.
Nervous system damage: Clinically significant structural neurological, or functional impairment
Popova et al. identified 428 ICD-10 conditions as co-occurring in individuals with FAS. Excluding conditions used in FAS diagnosis, co-occurring conditions with 50% prevalence or greater include:
Conduct disorder, behavioral problems, disruptive behavior, or impulsivity
Developmental language disorder
Hearing loss
Visual impairment, including blindness or astigmatism
Developmental delay, cognitive disorder, or mental deficiency
Premature birth
Substance dependence
Congenital malformation of retina
Congenital fusion of cervical vertebrae or cervical spine fusion
Attention deficit hyperactivity disorder or otherwise impaired attention
Small eye openings (blepharophimosis), or an abnormally increased distance between the eyes, or both (hypertelorism)
Other FASD conditions are partial expressions of FAS where the central nervous system shows clinical deficits. In these other FASD conditions, an individual may be at greater risk for adverse outcomes because brain damage is present without associated visual cues of poor growth or the "FAS face" that might ordinarily trigger an FASD evaluation. Such individuals may be misdiagnosed with primary mental health disorders such as ADHD or oppositional defiance disorder without appreciation that brain damage is the underlying cause of these disorders, which requires a different treatment paradigm than typical mental health disorders. While other FASD conditions may not yet be included as an ICD or DSM-IV-TR diagnosis, they nonetheless pose significant impairment in functional behavior because of underlying brain damage. Many indications of fetal alcohol spectrum disorders are developmental. Therefore, although a child may appear 'normal' at birth, intellectual disabilities caused by alcohol before birth may not appear until the child begins school.
More broadly, alcohol use during pregnancy is also associated with:
Intellectual disability, both in overall IQ measurements and in many functional tests
Fetal mortality, such as spontaneous abortion (Miscarriage), stillbirth, and sudden infant death syndrome
Domestic violence and potential harm to the infant
Sleep and sucking problems as a baby
Heart: A heart murmur that frequently disappears by one year of age. Ventricular septal defect, atrial septal defect, tetralogy of Fallot, coarctation of the aorta, or cardiac rhythm dysfunction.
Bones: Joint anomalies including abnormal position and function, altered palmar crease patterns, small distal phalanges, and small fifth fingernails.
Kidneys: Horseshoe, aplastic, dysplastic, or hypoplastic kidneys.
Eyes: Strabismus, optic nerve hypoplasia (which may cause light sensitivity, decreased visual acuity, or involuntary eye movements).
Cleft lip with or without a cleft palate: Alcohol is known to be a folic acid antagonist, and a baby's palate and lip develop during the first trimester of the pregnancy (first 12 weeks). Heavy alcohol consumption and binge drinking during this time have been linked to orofacial cleft.
Occasional problems: ptosis of the eyelid, microphthalmia, webbed neck, short neck, radioulnar synostosis, spina bifida, and hydrocephalus.
Alcohol can also harm the fertility of women who are planning for pregnancy. Adverse effects of alcohol can lead to malnutrition, seizures, vomiting and dehydration. The mother can suffer from anxiety and depression which can result in child abuse/neglect. It has also been observed that when the pregnant mother withdraws from alcohol, its effects are visible on the infant as well. The baby remains in an irritated mood, cries frequently, does not sleep properly, weakening of sucking ability and increased hunger.
In 2019, a study found that individuals with FASD have a higher risk of hypertension independent of race/ethnicity and obesity.
Causes
Fetal alcohol spectrum disorders are caused by alcohol exposure during gestational development. If an individual was not exposed to alcohol before birth, they will not have FASD. However, not all infants exposed to alcohol in utero will have detectable FAS, FASD, or pregnancy complications.
Exposure limit
No safe level of fetal alcohol exposure has been established. Because alcohol is a known teratogen, it is considered unethical to do randomized controlled trials on pregnant women to determine the precise toxicity effects of alcohol. Among women who consume any quantity of alcohol during pregnancy, the risk of giving birth to a child with FASD is about 15%, and to a child with FAS about 1.5%. Drinking 2 standard drinks a day, or 6 standard drinks in a short time, carries a 4.3% risk of a FAS birth (i.e. one of every 23 heavy-drinking pregnant women will deliver a child with FAS). Furthermore, alcohol-related congenital abnormalities occur at an incidence of roughly one out of 67 women who drink alcohol during pregnancy. Among those mothers who have an alcohol use disorder, an estimated one-third of their children have FAS. The variance seen in outcomes of alcohol consumption during pregnancy is poorly understood. Aggravating factors may include advanced maternal age, smoking, poor diet, genetics, and social risk factors.
The risk of FASD increases with amount consumed, the frequency of consumption, and longer duration of alcohol consumption during pregnancy. Blood alcohol concentration has been identified as a relevant factor. All forms of alcohol, such as beer, wine, and liquor, pose similar risk. Binge drinking increases the chances and severity of FASD to such an extent that Svetlana Popova has stated that "binge drinking is the direct cause of FAS or FASD". Small amounts of alcohol may not cause an abnormal appearance, however, small amounts of alcohol consumption while pregnant may cause behavioral problems and also increases the risk of miscarriage. Quasi-experimental studies provide moderately strong evidence that prenatal alcohol exposure causes detrimental cognitive outcomes, and some evidence of reduced birthweight, although no study was fully rated at low risk of bias and quantity of studies was limited.
The evidence is inconsistent and contradictory regarding the effects of low-to-moderate drinking, for example less than 12 grams of ethanol per day. Many studies find no significant effect, but some find beneficial associations, and others find detrimental associations, even on the same outcomes. Summarizing studies by country shows some similarity in results, due to using the same data sources. The definition of low alcohol consumption varies significantly among studies and often fails to incorporate all aspects of timing, dose, and duration. Recall bias and socioeconomic and psychosocial factors have been controlled for in most studies, but it is likely that residual confounding due to missing factors and variation in methods still exists and is larger than any observed effects.
Paternal Alcohol Use
Fathers who consume alcohol prior to conception may contribute to FAS through long term epigenetic modification of the father's sperm.
Prevention and Stigma
Almost all experts recommend that the mother abstain from alcohol use during pregnancy to prevent FASDs. A pregnant woman may not become aware that she has conceived until several weeks into the pregnancy, so it is also recommended to abstain from alcohol while attempting to become pregnant. The recommendations of abstaining from alcohol during pregnancy and while attempting to become have been made by the Surgeon General of the United States, the Centers for Disease Control, the American College of Obstetricians and Gynecologists, the American Academy of Pediatrics, the World Health Organization, the United Kingdom's National Institute for Health and Clinical Excellence, and many others. In the United States, federal legislation has required that warning labels be placed on all alcoholic beverage containers since 1988 under the Alcoholic Beverage Labeling Act.
Stigma
The most current advocacy perspectives encourage people and systems to approach FASD with interventions and support for individuals who are already living with FASD. Focusing on prevention often only further stigmatizes individuals with FASD and their birth parents. Advocates say, if a person is supporting people currently living with FASD then that person is spreading the awareness needed for successful prevention efforts; "Intervention is Prevention". Many social determinants of health impact the effects of PAE:
Genetics
Poverty/Access to nutritious food
Malnutrition
Poor social support networks
Lack of personal autonomy
Access to healthcare
Generational and sociocultural traumas
Access to mental health care and treatment
Medication
Women can experience serious symptoms that accompany alcohol withdrawal during pregnancy. According to the World Health Organization, these symptoms can be treated during pregnancy with brief use of benzodiazepine tranquilizers. Currently, the FDA has approved three medications—naltrexone, acamprosate, and disulfiram—for the treatment of alcohol use disorder (AUD). However, there is insufficient data regarding the safety of these medications for pregnant women.
Naltrexone is a nonselective opioid antagonist that is used to treat AUD and opioid use disorder. The long-term effects of naltrexone on the fetus are currently unknown. Animal studies show that naltrexone administered during pregnancy increases the incidence of early fetal loss; however, there are insufficient data available to identify the extent to which this is a risk in pregnant women.
Acamprosate functions as both an antagonist of NMDA and glutamate and an agonist at GABAA receptors, although its molecular mechanism is not completely understood. Acamprosate has been shown to be effective at preventing alcohol relapse during abstinence. Animal data, however, suggests that acamprosate can have possible teratogenic effects on fetuses.
Disulfiram prevents relapse by blocking the metabolism of acetaldehyde after consumption of alcohol which leads to headache, nausea, and vomiting. Some evidence suggests that disulfiram use during the first trimester is associated with an increased risk of congenital malformations such as reduction defects and cleft palate. Additionally, the effects of disulfiram can involve hypertension which can be harmful to both the pregnant woman and the fetus.
American Psychiatric Association guidelines recommend that medications not be used to treat alcohol use disorder in pregnant women except in cases of acute alcohol withdrawals or other co-existing conditions. Instead, behavioral interventions are usually preferred as treatments for pregnant women with AUD. Medications should only be used for pregnant women after carefully considering potential risks and harms of the medications versus the benefits of alcohol cessation.
Mechanism
After a pregnant woman consumes alcohol, the alcohol crosses through the placenta and umbilical cord to the developing fetus. Alcohol metabolizes slowly in the fetus and remains for a long time when compared to an adult. A human fetus appears to be at triple risk from maternal alcohol consumption:
The placenta allows free entry of ethanol and toxic metabolites like acetaldehyde into the fetal compartment. The so-called placental barrier is practically absent with respect to ethanol.
The developing fetal nervous system appears particularly sensitive to ethanol toxicity. The latter interferes with proliferation, differentiation, neuronal migration, axonic outgrowth, integration, and fine-tuning of the synaptic network. In short, all major processes in the developing central nervous system appear compromised.
Fetal tissues are quite different from adult tissues in function and purpose. For example, the main detoxicating organ in adults is the liver, whereas the fetal liver is incapable of detoxifying ethanol, as the ADH and ALDH enzymes have not yet been brought to expression at this early stage. Up to term, fetal tissues do not have significant capacity for the detoxification of ethanol, and the fetus remains exposed to ethanol in the amniotic fluid for periods far longer than the decay time of ethanol in the maternal circulation. The lack of significant quantities of ADH and ALDH means that fetal tissues have much lower quantities of antioxidant enzymes, like SOD, glutathione transferases, and glutathione peroxidases, resulting in antioxidant protection being much less effective.
Although alcohol is known to be a teratogen (causing birth defects), the exact biological mechanisms for the development of FAS or FASD are unknown. However, clinical and animal studies have identified a broad spectrum of pathways through which maternal alcohol can negatively affect the outcome of a pregnancy. Clear conclusions with universal validity are difficult to draw, since different ethnic groups show considerable genetic polymorphism for the hepatic enzymes responsible for ethanol detoxification. Genetic examinations have revealed a continuum of long-lasting molecular effects that are not only timing specific but are also dosage specific; with even moderate amounts being able to cause alterations. Additionally, ethanol may alter fetal development by interfering with retinoic acid signaling as acetaldehyde can compete with retinaldehyde and prevents its oxidation to retinoic acid.
Developmental stages
Different body systems in the infant grow, mature and develop at specific times during gestation. The effect of consumption of alcohol differs during each of these developmental stages:
From conception and to the third week, the most susceptible systems and organs are the brain, spinal cord, and heart. The effects of alcohol consumption early in the pregnancy can result in defects to these systems and organs.
During the third week, alcohol can also damage the central nervous system of the fetus.
During the fourth week of gestation, the limbs are being formed and it is at this point that alcohol can affect the development of arms, legs, fingers and toes. The eyes and ears also form during the fourth week and are more susceptible to the effects of alcohol.
By the sixth week of gestation, the teeth and palate are forming and alcohol consumption at this time will affect these structures. Alcohol use in this window is responsible for many of the facial characteristics of fetal alcohol syndrome.
During the twelfth week, frequent alcohol exposure can negatively impact the brain development which affects cognitive, learning and behavioral skills before birth.
By the 20th week of gestation the formation of organs and organ systems is well-developed. The infant is still susceptible to the damaging effects of alcohol.
Ethanol exposure in the second trimester reduces nutrition levels and can affect the functioning of the endocrine system in both fetus and mother. This is because blood flow via umbilical artery to fetal brain is reduced.
Diagnosis
Fetal alcohol spectrum disorders encompass a range of physical and neurodevelopmental problems which can result from prenatal alcohol exposure. Diagnosis is based on the signs and symptoms in the person and evidence of alcohol use. These diagnoses of fetal alcohol spectrum disorders are currently recognized:
Fetal alcohol syndrome (FAS)
Partial fetal alcohol syndrome (pFAS) refers to individuals with a known, or highly suspected, history of prenatal alcohol exposure who have alcohol-related physical and neurodevelopmental deficits that do not meet the full criteria for FAS.
Alcohol-related neurodevelopmental disorder (ARND)
Neurobehavioral disorder associated with prenatal alcohol exposure (ND-PAE)
As of 2016, the Swedish Agency for Health Technology Assessment and Assessment of Social Services accepts only FAS as a diagnosis, seeing the evidence as inconclusive with respect to other types. The agency feels it is unclear if identifying a FASD-related condition benefits the diagnosed individual.
Classification
Presently, four FASD diagnostic systems that diagnose FAS and other FASD conditions have been developed in North America:
The Institute of Medicine's guidelines for FAS, the first system to standardize diagnoses of individuals with prenatal alcohol exposure;
The University of Washington's "The 4-Digit Diagnostic Code", which ranks the four key features of FASD on a Likert scale of one to four and yields 256 descriptive codes that can be categorized into 22 distinct clinical categories, ranging from FAS to no findings;
The Centers for Disease Control's "Fetal Alcohol Syndrome: Guidelines for Referral and Diagnosis", which established consensus on the diagnosis FAS in the U.S. but deferred addressing other FASD conditions; and
Canadian guidelines for FASD diagnoses, which established criteria for diagnosing FASD in Canada and harmonized most differences between the IOM and University of Washington's systems.
Each diagnostic system requires an assessment of four key features: growth, facial features, central nervous system, and alcohol exposure. To determine any FASD condition, a multi-disciplinary evaluation is necessary to assess each of the four key features for assessment. Generally, a trained physician will determine growth deficiency and FAS facial features. While a qualified physician may also assess central nervous system structural abnormalities or neurological problems, usually central nervous system damage is determined through psychological, speech-language, and occupational therapy assessments to ascertain clinically significant impairments in three or more of the Ten Brain Domains. Prenatal alcohol exposure risk may be assessed by a qualified physician, psychologist, social worker, or chemical health counselor. These professionals work together as a team to assess and interpret data of each key feature for assessment and develop an integrative, multi-disciplinary report to diagnose FAS (or other FASD conditions) in an individual.
A positive finding on all four features is required for a diagnosis of FAS, and the four diagnostic systems essentially agree on criteria for fetal alcohol syndrome (FAS). However, there are differences among systems when full criteria for FAS are not met. Prenatal alcohol exposure and central nervous system damage are the critical elements of the spectrum of FASD, and a positive finding in these two features is sufficient for an FASD diagnosis in all FASD systems. But different researchers and systems may use a wide variety of terminology to describe an individual's FASD condition, as the nomenclature is still evolving. Most individuals with deficits resulting from prenatal alcohol exposure do not express all features of FAS and fall into other FASD conditions. The Canadian guidelines recommend the assessment and descriptive approach of the "4-Digit Diagnostic Code" for each key feature of FASD and the terminology of the IOM in diagnostic categories, excepting ARBD.
Fetal alcohol syndrome
The most severe condition is called Fetal Alcohol Syndrome (FAS), which refers to individuals who have a specific set of birth defects and neurodevelopmental disorders characteristic of the diagnosis. The following criteria must be fully met for an FAS diagnosis:
Prenatal or postnatal height or weight (or both) at or below the 10th percentile
All three FAS facial features present
Clinically significant structural, neurological, or functional impairment of the central nervous system
Confirmed or Unknown prenatal alcohol exposure
FAS is the only expression of FASD that has garnered consensus among experts to become an official ICD-9 and ICD-10 diagnosis.
Partial FAS
Partial FAS (pFAS) was previously known as atypical FAS in the 1997 edition of the "4-Digit Diagnostic Code". People with pFAS have a confirmed history of prenatal alcohol exposure, but may lack growth deficiency or the complete facial stigmata. Central nervous system damage is present at the same level as FAS. These individuals have the same functional disabilities but "look" less like FAS.
The following criteria must be fully met for a diagnosis of Partial FAS:
Two or three FAS facial features present
Clinically significant structural, neurological, or functional impairment in three or more of the Ten Brain Domains
Confirmed prenatal alcohol exposure
Growth or height may range from normal to deficient.
Alcohol-related neurodevelopmental disorder
Alcohol-related neurodevelopmental disorder (ARND) is the specific diagnosis of the non-dysmorphic type of FASD, where a majority of the symptoms are witnessed. The diagnosis was initially suggested by the Institute of Medicine to replace the terms FAE (fetal alcohol effects). It focuses on central nervous system damage, rather than growth deficiency or FAS facial features. The Canadian guidelines also use this diagnosis and the same criteria. While the "4-Digit Diagnostic Code" includes these criteria for three of its diagnostic categories, it refers to this condition as static encephalopathy. The behavioral effects of ARND are not necessarily unique to alcohol however, so use of the term must be within the context of confirmed prenatal alcohol exposure. ARND may be gaining acceptance over the terms FAE and ARBD to describe FASD conditions with central nervous system abnormalities or behavioral or cognitive abnormalities or both due to prenatal alcohol exposure without regard to growth deficiency or FAS facial features.
The following criteria must be fully met for a diagnosis of ARND or static encephalopathy:
Minimal or no FAS facial features present
Clinically significant structural, neurological, or functional impairment in three or more of the Ten Brain Domains
Confirmed prenatal alcohol exposure
Growth or height may range from normal to minimally deficient.
Neurobehavioral disorder associated with prenatal alcohol exposure
Neurobehavioral disorder associated with prenatal alcohol exposure (ND-PAE) is the spectrum-wide term for the psychiatric, behavioral, and neurological symptoms of all FASDs. It was introduced into the DSM-V as a "condition for further study" and as a specified condition under, "other specified neurodevelopmental disorders" as a way to better study the behavioral aspects of all FASD disorders.
Specific criteria
Growth
In terms of FASD, growth deficiency is defined as significantly below average height, weight or both due to prenatal alcohol exposure and can be assessed at any point in the lifespan. Growth measurements must be adjusted for parental height, gestational age (for a premature infant), and other postnatal insults (e.g., poor nutrition), although birth height and weight are the preferred measurements. Deficiencies are documented when height or weight falls at or below the 10th percentile of standardized growth charts appropriate to the population. Prenatal or postnatal presentation of growth deficits can occur, but are most often postnatal.
Criteria for FASD are least specific in the Institute of Medicine (IOM) diagnostic system ("low birth weight..., decelerating weight not due to nutrition..., [or] disproportional low weight to height" p. 4 of executive summary), while the CDC use the 10th percentile as a cut-off to determine growth deficiency. The "4-Digit Diagnostic Code" allows for mid-range gradations in growth deficiency (between the 3rd and 10th percentiles) and severe growth deficiency at or below the 3rd percentile. Growth deficiency (at severe, moderate, or mild levels) contributes to diagnoses of FAS and pFAS, but not ARND or static encephalopathy.
The "4-Digit Diagnostic Code" from 2004 ranks growth deficiency as follows:
Severe: Height and weight at or below the 3rd percentile.
Moderate: Either height or weight at or below the 3rd percentile, but not both.
Mild: Either height or weight or both between the 3rd and 10th percentiles.
None: Height and weight both above the 10th percentile.
In the initial studies that described FAS, growth deficiency was a requirement for inclusion in the studies; thus, all the original people with FAS had growth deficiency as an artifact of sampling characteristics used to establish criteria for the syndrome. That is, growth deficiency is a key feature of FASD because growth deficiency was a criterion for inclusion in the study that defined FAS. Growth deficiency may be less critical for understanding the disabilities of FASD than the neurobehavioral sequelae to the brain damage. Canadian guidelines updated in 2016 deleted growth as a diagnostic criterion.
Facial features
Several characteristic craniofacial abnormalities are often visible in individuals with FAS. The presence of FAS facial features indicates brain damage, although brain damage may also exist in their absence. FAS facial features (and most other visible, but non-diagnostic, deformities) are believed to be caused mainly during the 10th to 20th week of gestation.
Refinements in diagnostic criteria since 1975 have yielded three distinctive and diagnostically significant facial features which distinguish FAS from other disorders with partially overlapping characteristics. The three FAS facial features are:
A smooth philtrum: The divot or groove between the nose and upper lip flattens with increased prenatal alcohol exposure.
Thin vermilion: The upper lip thins with increased prenatal alcohol exposure.
Small palpebral fissures: Eye width decreases with increased prenatal alcohol exposure.
Measurement of FAS facial features uses criteria developed by the University of Washington. The lip and philtrum are measured by a trained physician with the Lip-Philtrum Guide, a five-point Likert scale with representative photographs of lip and philtrum combinations ranging from normal (ranked 1) to severe (ranked 5). Palpebral fissure length (PFL) is measured in millimeters with either calipers or a clear ruler and then compared to a PFL growth chart, also developed by the University of Washington.
Ranking FAS facial features is complicated because the three separate facial features can be affected independently by prenatal alcohol. A summary of the criteria follows:
Severe: All three facial features ranked independently as severe (lip ranked at 4 or 5, philtrum ranked at 4 or 5, and PFL two or more standard deviations below average).
Moderate: Two facial features ranked as severe and one feature ranked as moderate (lip or philtrum ranked at 3, or PFL between one and two standard deviations below average).
Mild: A mild ranking of FAS facial features covers a broad range of facial feature combinations:
Two facial features ranked severe and one ranked within normal limits,
One facial feature ranked severe and two ranked moderate, or
One facial feature ranked severe, one ranked moderate and one ranked within normal limits.
None: All three facial features ranked within normal limits.
Central nervous system
Central nervous system (CNS) damage is the primary feature of any FASD diagnosis. Prenatal alcohol exposure, which is classified as a teratogen, can damage the brain across a continuum of gross to subtle impairments, depending on the amount, timing, and frequency of the exposure as well as genetic predispositions of the fetus and mother. While functional abnormalities are the behavioral and cognitive expressions of the FASD disability, CNS damage can be assessed in three areas: structural, neurological, and functional impairments.
All four diagnostic systems allow for assessment of CNS damage in these areas, but criteria vary. The IOM system requires structural or neurological impairment for a diagnosis of FAS, but also allows a "complex pattern" of functional anomalies for diagnosing PFAS and ARND. The "4-Digit Diagnostic Code" and CDC guidelines allow for a positive CNS finding in any of the three areas for any FASD diagnosis, but functional anomalies must measure at two standard deviations or worse in three or more functional domains for a diagnosis of FAS, PFAS, and ARND. The "4-Digit Diagnostic Code" also allows for an FASD diagnosis when only two functional domains are measured at two standard deviations or worse. The "4-Digit Diagnostic Code" further elaborates the degree of CNS damage according to four ranks:
Definite: Structural impairments or neurological impairments for FAS or static encephalopathy.
Probable: Significant dysfunction of two standard deviations or worse in three or more functional domains.
Possible: Mild to moderate dysfunction of two standard deviations or worse in one or two functional domains or by judgment of the clinical evaluation team that CNS damage cannot be dismissed.
Unlikely: No evidence of CNS damage.
Structural
Structural abnormalities of the brain are observable, physical damage to the brain or brain structures caused by prenatal alcohol exposure. Structural impairments may include microcephaly (small head size) of two or more standard deviations below the average, or other abnormalities in brain structure (e.g., agenesis of the corpus callosum, cerebellar hypoplasia).
Microcephaly is determined by comparing head circumference (often called occipitofrontal circumference, or OFC) to appropriate OFC growth charts. Other structural impairments must be observed through medical imaging techniques by a trained physician. Because imaging procedures are expensive and relatively inaccessible to most people, diagnosis of FAS is not frequently made via structural impairments, except for microcephaly.
Evidence of a CNS structural impairment due to prenatal alcohol exposure will result in a diagnosis of FAS, and neurological and functional impairments are highly likely.
During the first trimester of pregnancy, alcohol interferes with the migration and organization of brain cells, which can create structural deformities or deficits within the brain. During the third trimester, damage can be caused to the hippocampus, which plays a role in memory, learning, emotion, and encoding visual and auditory information, all of which can create neurological and functional CNS impairments as well.
As of 2002, there were 25 reports of autopsies on infants known to have FAS. The first was in 1973, on an infant who died shortly after birth. The examination revealed extensive brain damage, including microcephaly, migration anomalies, corpus callosum dysgenesis, and a massive neuroglial, leptomeningeal heterotopia covering the left hemisphere.
In 1977, Clarren described a second infant whose mother was a binge drinker. The infant died ten days after birth. The autopsy showed severe hydrocephalus, abnormal neuronal migration, and a small corpus callosum. FAS has also been linked to brainstem and cerebellar changes, agenesis of the corpus callosum and anterior commissure, neuronal migration errors, absent olfactory bulbs, meningomyelocele, and porencephaly.
Neurological
When structural impairments are not observable or do not exist, neurological impairments are assessed. In the context of FASD, neurological impairments are caused by prenatal alcohol exposure which causes general neurological damage to the central nervous system (CNS), the peripheral nervous system, or the autonomic nervous system. A determination of a neurological problem must be made by a trained physician, and must not be due to a postnatal insult, such as meningitis, concussion, traumatic brain injury, etc.
All four diagnostic systems show virtual agreement on their criteria for CNS damage at the neurological level, and evidence of a CNS neurological impairment due to prenatal alcohol exposure will result in a diagnosis of FAS or pFAS, and functional impairments are highly likely.
Neurological problems are expressed as either hard signs, or diagnosable disorders, such as epilepsy or other seizure disorders, or soft signs. Soft signs are broader, nonspecific neurological impairments, or symptoms, such as impaired fine motor skills, neurosensory hearing loss, poor gait, clumsiness, and poor eye–hand coordination. Many soft signs have norm-referenced criteria, while others are determined through clinical judgment. "Clinical judgment" is only as good as the clinician, and soft signs should be assessed by either a pediatric neurologist, a pediatric neuropsychologist, or both.
Functional
When structural or neurological impairments are not observed, all four diagnostic systems allow CNS damage due to prenatal alcohol exposure to be assessed in terms of functional impairments. Functional impairments are deficits, problems, delays, or abnormalities due to prenatal alcohol exposure (rather than hereditary causes or postnatal insults) in observable and measurable domains related to daily functioning, often referred to as developmental disabilities. There is no consensus on a specific pattern of functional impairments due to prenatal alcohol exposure and only CDC guidelines label developmental delays as such, so criteria (and FASD diagnoses) vary somewhat across diagnostic systems.
The four diagnostic systems list various CNS domains that can qualify for functional impairment that can determine an FASD diagnosis:
Evidence of a complex pattern of behavior or cognitive abnormalities inconsistent with developmental level in the following CNS domains – Sufficient for a pFAS or ARND diagnosis using IOM guidelines
Learning disabilities, academic achievement, impulse control, social perception, communication, abstraction, math skills, memory, attention, judgment
Performance at two or more standard deviations on standardized testing in three or more of the following CNS domains – Sufficient for an FAS, pFAS or static encephalopathy diagnosis using 4-Digit Diagnostic Code
Executive functioning, memory, cognition, social/adaptive skills, academic achievement, language, motor skills, attention, activity level
General cognitive deficits (e.g., IQ) at or below the 3rd percentile on standardized testing – Sufficient for an FAS diagnosis using CDC guidelines
Performance at or below the 16th percentile on standardized testing in three or more of the following CNS domains – Sufficient for an FAS diagnosis using CDC guidelines
Cognition, executive functioning, motor functioning, attention and hyperactive problems, social skills, sensory processing disorder, social communication, memory, difficulties responding to common parenting practices
Performance at two or more standard deviations on standardized testing in three or more of the following CNS domains – Sufficient for an FAS diagnosis using Canadian guidelines
Cognition, communication, academic achievement, memory, executive functioning, adaptive behavior, motor skills, social skills, social communication
Ten brain domains
A recent effort to standardize assessment of functional CNS damage has been suggested by an experienced FASD diagnostic team in Minnesota. The proposed framework attempts to harmonize IOM, 4-Digit Diagnostic Code, CDC, and Canadian guidelines for measuring CNS damage vis-à-vis FASD evaluations and diagnosis. The standardized approach is referred to as the Ten Brain Domains and encompasses aspects of all four diagnostic systems' recommendations for assessing CNS damage due to prenatal alcohol exposure. The framework provides clear definitions of brain dysfunction, specifies empirical data needed for accurate diagnosis, and defines intervention considerations that address the complex nature of FASD with the intention to avoid common secondary disabilities.
The proposed Ten Brain Domains include:
Achievement
Adaptive behavior
Attention
Cognition
Executive functioning
Language
Memory
Motor skills
Multisensory integration or soft neurological problems
Social communication
The Fetal Alcohol Diagnostic Program (FADP) uses unpublished Minnesota state criteria of performance at 1.5 or more standard deviations on standardized testing in three or more of the Ten Brain Domains to determine CNS damage. However, the Ten Brain Domains are easily incorporated into any of the four diagnostic systems' CNS damage criteria, as the framework only proposes the domains, rather than the cut-off criteria for FASD.
Alcohol exposure
Prenatal alcohol exposure is determined by interview of the biological mother or other family members knowledgeable of the mother's alcohol use during the pregnancy (if available), prenatal health records (if available), and review of available birth records, court records (if applicable), chemical dependency treatment records (if applicable), chemical biomarkers, or other reliable sources.
Exposure level is assessed as confirmed exposure, unknown exposure, and confirmed absence of exposure by the IOM, CDC and Canadian diagnostic systems. The "4-Digit Diagnostic Code" further distinguishes confirmed exposure as High Risk and Some Risk:
Confirmed exposure: The CDC guidelines are silent on using information on amount, frequency, and timing of prenatal alcohol use for diagnostic purposes. The IOM and Canadian guidelines explore this further, acknowledging the importance of significant alcohol exposure from regular or heavy episodic alcohol consumption in determining, but offer no standard for diagnosis. Canadian guidelines discuss this lack of clarity and parenthetically point out that "heavy alcohol use" is defined by the National Institute on Alcohol Abuse and Alcoholism as five or more drinks per episode on five or more days during a 30-day period. "The 4-Digit Diagnostic Code" ranking system distinguishes between levels of prenatal alcohol exposure as high risk and some risk. It operationalizes high risk exposure as a blood alcohol concentration (BAC) greater than 100 mg/dL delivered at least weekly in early pregnancy. This BAC level is typically reached by a 55 kg female drinking six to eight beers in one sitting.
High Risk: Confirmed use of alcohol during pregnancy known to be at high blood alcohol levels (100 mg/dL or greater) delivered at least weekly in early pregnancy.
Some Risk: Confirmed use of alcohol during pregnancy with use less than High Risk or unknown usage patterns.
Unknown Risk: Unknown use of alcohol during pregnancy. For many adopted or adults and children in foster care, records or other reliable sources may not be available for review. Reporting alcohol use during pregnancy can also be stigmatizing to birth mothers, especially if alcohol use is ongoing. Many are reluctant to admit to drinking or to provide an accurate report of the quantity they drank. In these cases, all diagnostic systems use an unknown prenatal alcohol exposure designation. A diagnosis of FAS is still possible with an unknown exposure level if other key features of FASD are present at clinical levels.
No Risk: Confirmed absence of prenatal alcohol exposure. Confirmed absence of exposure would apply to planned pregnancies in which no alcohol was used or pregnancies of women who do not use alcohol or report no use during the pregnancy. This designation is relatively rare, as most people presenting for an FASD evaluation are at least suspected to have had a prenatal alcohol exposure due to presence of other key features of FASD.
Biomarkers
Evidence is insufficient for the use of chemical biomarkers to detect prenatal alcohol exposure. Biomarkers being studied include fatty acid ethyl esters (FAEE) detected in the meconium (first feces of an infant) and hair. FAEE may be present if chronic alcohol exposure occurs during the second and third trimester since this is when the meconium begins to form. Concentrations of FAEE can be influence by medication use, diet, and individual genetic variations in FAEE metabolism however.
Differential diagnosis
The CDC reviewed nine syndromes that have overlapping features with FAS; however, none of these syndromes include all three FAS facial features, and none are the result of prenatal alcohol exposure:
Aarskog syndrome
Williams syndrome
Noonan syndrome
Dubowitz syndrome
Brachman-DeLange syndrome
Toluene syndrome
Fetal hydantoin syndrome
Fetal valproate syndrome
Maternal PKU fetal effects
Other disorders that have overlapping behavioral symptoms that might be comorbid to fetal alcohol spectrum disorder might include:
Attention deficit hyperactive disorder
Autism spectrum disorder
Reactive attachment disorder
Oppositional defiant disorder
Sensory integration dysfunction
Bipolar disorder
Depression
Most people with an FASD have most often been misdiagnosed with ADHD due to the large overlap between their behavioral deficits.
Treatment
Although the condition has no available cure, treatment can improve outcomes. Because CNS damage, symptoms, secondary disabilities, and needs vary widely by individual, there is no one treatment type that works for everyone.
Between 2017 and 2019 researchers made a breakthrough when they discovered a possible cure using neural stem cells (NSCs); they propose that if applied to a newborn, the damage can be reversed and prevent any lasting effects in the future.
Medications
Psychoactive drugs are frequently tried as many FASD symptoms are mistaken for or overlap with other disorders, most notably ADHD. Medications are used to specifically treat symptoms of FASDs and not FAS entirely. Some of the medications used are antidepressants, stimulants, neuroleptics and anti-anxiety drugs.
Behavioral interventions
Early intervention from birth to age 3 has been shown to improve the development of a child born with FASD. Interventions may include parent–child interaction therapy, efforts to modify child behavior, and drugs. Behavioral interventions are based on the learning theory, which is the basis for many parenting and professional strategies and interventions. Along with ordinary parenting styles, such strategies are frequently used by default for treating those with FAS, as the diagnoses oppositional defiance disorder (ODD), conduct disorder, reactive attachment disorder (RAD) often overlap with FAS (along with ADHD), and these are sometimes thought to benefit from behavioral interventions. Frequently, a person's poor academic achievement results in special education services, which also utilizes principles of learning theory, behavior modification, and outcome-based education.
Children with FAS benefit from behavioral and functional training, social skill training and tutoring. Support groups and talk therapy not only help the children suffering from FAS, but also help the parents and siblings of these children.
Developmental framework
Many books and handouts on FAS recommend a developmental approach, based on developmental psychology, even though most do not specify it as such and provide little theoretical background. Optimal human development generally occurs in identifiable stages (e.g., Jean Piaget's theory of cognitive development, Erik Erikson's stages of psychosocial development, John Bowlby's attachment framework, and other developmental stage theories). FAS interferes with normal development, which may cause stages to be delayed, skipped, or immaturely developed. Over time, an unaffected child can negotiate the increasing demands of life by progressing through stages of development normally, but not so for a child with FAS.
By knowing what developmental stages and tasks children follow, treatment and interventions for FAS can be tailored to helping a person meet developmental tasks and demands successfully. If a person is delayed in the adaptive behavior domain, for instance, then interventions would be recommended to target specific delays through additional education and practice (e.g., practiced instruction on tying shoelaces), giving reminders, or making accommodations (e.g., using slip-on shoes) to support the desired functioning level. This approach is an advance over behavioral interventions, because it takes the person's developmental context into account while developing interventions.
Advocacy model
The advocacy model takes the point of view that someone is needed to actively mediate between the environment and the person with FAS. Advocacy activities are conducted by an advocate (for example, a family member, friend, or case manager) and fall into three basic categories. An advocate for FAS: (1) interprets FAS and the disabilities that arise from it and explains it to the environment in which the person operates, (2) engenders change or accommodation on behalf of the person, and (3) assists the person in developing and reaching attainable goals.
The advocacy model is often recommended, for example, when developing an individualized education program (IEP) for the person's progress at school.
An understanding of the developmental framework would presumably inform and enhance the advocacy model, but advocacy also implies interventions at a systems level as well, such as educating schools, social workers, and so forth on best practices for FAS. However, several organizations devoted to FAS also use the advocacy model at a community practice level as well.
Treating FAS at the public health and public policy level promotes FAS prevention and diversion of public resources to assist those with FAS. It is related to the advocacy model but promoted at a systems level (rather than with the individual or family), such as developing community education and supports, state or province level prevention efforts (e.g., screening for maternal alcohol use during OB/GYN or prenatal medical care visits), or national awareness programs. Several organizations and state agencies in the U.S. are dedicated to this type of intervention.
Prognosis
The prognosis of FASD is variable depending on the type, severity, and if treatment is issued. Prognostic disabilities are divided into primary and secondary disabilities.
Primary disabilities
The primary disabilities of FAS are the functional difficulties with which the child is born as a result of CNS damage due to prenatal alcohol exposure.
Often, primary disabilities are mistaken as behavior problems, but the underlying CNS damage is the originating source of a functional difficulty, rather than a mental health condition, which is considered a secondary disability. The exact mechanisms for functional problems of primary disabilities are not always fully understood, but animal studies have begun to shed light on some correlates between functional problems and brain structures damaged by prenatal alcohol exposure. Representative examples include:
Learning impairments are associated with impaired dendrites of the hippocampus
Impaired motor development and functioning are associated with reduced size of the cerebellum
Hyperactivity is associated with decreased size of the corpus callosum
Functional difficulties may result from CNS damage in more than one domain, but common functional difficulties by domain include: Note that this is not an exhaustive list of difficulties.
Achievement: Learning disabilities
Adaptive behavior: Poor impulse control, poor personal boundaries, poor anger management, stubbornness, intrusive behavior, too friendly with strangers, poor daily living skills, developmental delays
Attention: Attention-deficit/hyperactivity disorder (ADHD), poor attention or concentration, distractible
Cognition: Intellectual disability, confusion under pressure, poor abstract skills, difficulty distinguishing between fantasy and reality, slower cognitive processing
Executive functioning: Poor judgment, information-processing disorder, poor at perceiving patterns, poor cause and effect reasoning, inconsistent at linking words to actions, poor generalization ability
Language: Expressive or receptive language disorders, grasp parts but not whole concepts, lack understanding of metaphor, idioms, or sarcasm
Memory: Poor short-term memory, inconsistent memory and knowledge base
Motor skills: Poor handwriting, poor fine motor skills, poor gross motor skills, delayed motor skill development (e.g., riding a bicycle at appropriate age)
Sensory processing and soft neurological problems: sensory processing disorder, sensory defensiveness, undersensitivity to stimulation
Social communication: Intrude into conversations, inability to read nonverbal or social cues, "chatty" but without substance
Secondary disabilities
The secondary disabilities of FAS are those that arise later in life secondary to CNS damage. These disabilities often emerge over time due to a mismatch between the primary disabilities and environmental expectations; secondary disabilities can be ameliorated with early interventions and appropriate supportive services.
Six main secondary disabilities were identified in a University of Washington research study of 473 subjects diagnosed with FAS, PFAS (partial fetal alcohol syndrome), and ARND (alcohol-related neurodevelopmental disorder):
Mental health problems: Diagnosed with ADHD, clinical depression, or other mental illness, experienced by over 90% of the subjects
Disrupted school experience: Suspended or expelled from school or dropped out of school, experienced by 60% of the subjects (age 12 and older)
Trouble with the law: Charged with or convicted of a crime, experienced by 60% of the subjects (age 12 and older)
Confinement: For inpatient psychiatric care, inpatient chemical dependency care, or incarcerated for a crime, experienced by about 50% of the subjects (age 12 and older)
Inappropriate sexual behavior: Sexual advances, sexual touching, or promiscuity, experienced by about 50% of the subjects (age 12 and older)
Alcohol and drug problems: Abuse or dependency, experienced by 35% of the subjects (age 12 and older)
Two additional secondary disabilities exist for adults:
Dependent living: Group home, living with family or friends, or some sort of assisted living, experienced by 80% of the subjects (age 21 and older)
Problems with employment: Required ongoing job training or coaching, could not keep a job, unemployed, experienced by 80% of the subjects (age 21 and older)
Protective factors and strengths
Eight factors were identified in the same study as universal protective factors that reduced the incidence rate of the secondary disabilities:
Living in a stable and nurturing home for over 73% of life
Being diagnosed with FAS before age six
Never having experienced violence
Remaining in each living situation for at least 2.8 years
Experiencing a "good quality home" (meeting 10 or more defined qualities) from age 8 to 12 years old
Having been found eligible for developmental disability (DD) services
Having basic needs met for at least 13% of life
Having a diagnosis of FAS (rather than another FASD condition)
Malbin (2002) has identified the following areas of interests and talents as strengths that often stand out for those with FASD and should be utilized, like any strength, in treatment planning:
Music, playing instruments, composing, singing, art, spelling, reading, computers, mechanics, woodworking, skilled vocations (welding, electrician, etc.), writing, poetry
Participation in non-impact sport or physical fitness activities
Lifespan
One study found that the people with FAS had a significantly shorter life expectancy. With the average life span of 34 years old, a study found that 44% of the deaths were of "external cause", with 15% of deaths being suicides.
Epidemiology
Globally, one in 10 women drink alcohol during pregnancy. Out of this population, 20% binge drink and have four or more alcoholic drinks per single occasion. The use of alcohol during pregnancy occurs at different rates across the world, potentially due to various cultural differences and legislation. The five countries with the highest prevalence of alcohol use during pregnancy are Ireland (60%), Belarus (47%), Denmark (46%), the UK (41%), and the Russian Federation (37%).
In a recent count, the prevalence of having any FASD disorder was 1 person out of 20, but some people estimate it could be as high as 1 in 7. The rates of FAS and FASD are likely to be underestimated, because of the difficulty in making the diagnosis and the reluctance of clinicians to label children and mothers.
Australia
FASD among Australian youth is more common in indigenous Australians. The only states that have registered birth defects in Australian youth are Western Australia, New South Wales, Victoria and South Australia. In Australia, only 12% of Australian health professionals are aware of the diagnostics and symptoms of FASD. In Western Australia, the rate of births resulting in FASD is 0.02 per 1,000 births for non-Indigenous Australians, however among indigenous births the rate is 2.76 per 1,000 births. In Victoria, there have been no registered FASD related births for indigenous Australians, but the rate for the general population in Victoria is 0.01–0.03 per 1000 births. There have been no dedicated FASD clinics within Western Australia, but there are also no nationally supported diagnostic criteria anywhere in Australia. Passive surveillance is a prevention technique used within Australia to assist in monitoring and establishing detectable defects during pregnancy and childhood.
Canada
A 2015 review article estimated the overall costs to Canada from FASD at $9.7 billion (including from crime, healthcare, education, etc.).
South Africa
In South Africa, some populations have rates as high as 9%.
United States
In the United States, alcohol use at some point during pregnancy is common and appears to be rising in prevalence. In 2006–2010, an estimated 7.6% of pregnant women used alcohol, while 1.4% of pregnant women reported binge drinking during their pregnancy. The highest prevalence estimates of reported alcohol use during pregnancy were among women who are aged 35–44 years (14.3%), white (8.3%), college graduates (10.0%), or employed (9.6%). In 2015, about 10% of pregnant women drank alcohol in the past month, and 20% to 30% drank at some point during the pregnancy. Of pregnant American women, 3.6% met criteria for an alcohol use disorder in a 2001 epidemiological study. As of 2016, the US Centers for Disease Control estimated 3 million women in the United States are at risk of having a baby with FASD.
FASD is estimated to affect between 1-2% and 5% of people in the United States and Western Europe. FAS is believed to occur in between 0.2 and 9 per 1,000 live births in the United States. Using medical and other records, CDC studies have identified 0.2 to 1.5 infants with FAS for every 1,000 live births in certain areas of the United States. A more recent CDC study of 2010 data analyzed medical and other records and found FAS in 0.3 out of 1,000 children from 7 to 9 years of age.
The lifetime cost per child with FAS in the United States was estimated at $2 million (for an overall cost across the country of over $4 billion) by the CDC in 2002.
History
Before designation
Some hold that ancient sources describe the negative effects of alcohol during pregnancy, identifying admonitions from ancient Greek, Roman, the Talmud, and the Bible. For example, Plato writes in his fourth-century B.C. Laws (6.775): "Drinking to excess is a practice that is nowhere seemly ... nor yet safe. ... It behooves both bride and bridegroom to be sober ... in order to ensure, as far as possible, in every case that the child that is begotten may be sprung from the loins of sober parents." The sixth-century AD Talmud (Kethuboth 60b) cautions, "One who drinks intoxicating liquor will have ungainly children." However, ancient sources rarely if ever distinguish maternal alcohol consumption from paternal, and are more concerned with conception than pregnancy. The sources can often be viewed as expressing heredity, that children are likely to turn out like their (alcoholic) parents, rather than presenting the modern viewpoint that alcohol itself has an impact.
In 1725, in the midst of the Gin Craze, British physicians petitioned the House of Commons on the effects of strong drink when consumed by pregnant women saying that such drinking is "too often the cause of weak, feeble, and distempered children, who must be, instead of an advantage and strength, a charge to their country". There are many other such historical references during that period. Gin specifically was implicated as affecting children's health and causing stillbirth and infant mortality, as depicted in William Hogarth's Gin Lane. In contrast, Hogarth's Beer Street shows commerce and happiness, suggesting that the alcohol in beer was not known to have deleterious effects at this time.
In Gaelic Scotland, according to Martin Martin, the mother and nurse were not allowed to consume ale during pregnancy and breastfeeding.
In the 19th century, Benjamin Rush and Thomas Trotter lobbied against alcohol consumption during pregnancy to avoid dependence and mental deficiency in children. The teetotalism and temperance movements popularized these and other claims, including the teratogenic effects of alcohol on animal embryos, but they were often sensationalized, so much so that any finding that alcohol was harmful was largely denounced as propaganda. A prominent observation of possible links between maternal alcohol use and fetal damage was made in 1899 by Dr. William Sullivan, a Liverpool prison physician who noted higher rates of stillbirth for 120 alcoholic female prisoners than their sober female relatives. He suggested the causal agent to be alcohol use. This contradicted the predominating belief at the time that heredity caused intellectual disability, poverty, and criminal behavior, which contemporary studies on the subjects usually concluded. A case study by Henry H. Goddard of the Kallikak family—popular in the early 1900s—represents this earlier perspective, though later researchers have suggested that the Kallikaks almost certainly had FAS. General studies and discussions on alcoholism throughout the mid-1900s were typically based on a heredity argument. Researchers were often temperance advocates and funded by like-minded organizations such as the Anti-Saloon League, so clinicians viewed all such research with heavy skepticism. The temperance movement effectively shut down serious research into the subject for nearly 50 years after Prohibition.
From the 1960s to the 1980s, alcohol was commonly used as a tocolytic, a method to stop preterm labor (born at less than 37 weeks gestation). The method originated with Dr. Fritz Fuchs, the chairman of the department of obstetrics and gynecology at Cornell University Medical College. Doctors recommended a small amount of alcohol to calm the uterus during contractions in early pregnancy or Braxton Hicks contractions. In later stages of pregnancy, the alcohol was administered intravenously and often in large amounts. "Women experienced similar effects as occur with oral ingestion, including intoxication, nausea and vomiting, and potential alcohol poisoning, followed by hangovers when the alcohol was discontinued." Vomiting put the mother at a high risk for aspiration and was "a brutal procedure for all involved". Because the alcohol was being given intravenously, the doctor could continue giving the treatment to the mother long after she had passed out, resulting in her being more intoxicated than would otherwise be possible. Such heavy intoxication was highly likely to contribute to FASD. In a 2015 review, ethanol was found to be no better than placebo (sugar water) in suppressing preterm birth and neonatal mortality. Not only was ethanol worse than other beta-mimetic drugs (tocolytic agents) at postponing birth, it also led to a higher rate of low birthweight babies, babies with breathing problems at birth, and neonatal death.
Recognition as a syndrome
In France in 1957, Jacqueline Rouquette had described 100 children whose parents were alcoholics in a thesis, which was not published. "She gave a good description in certain cases of the facies" according to her mentor, Paul Lemoine.
In 1968, Paul Lemoine of Nantes, himself published a study in a French medical journal about children with distinctive features whose mothers were alcoholics.
Independently, in the U.S., Christy Ulleland at University of Washington Medical School conducted an 18-month study in 1968–1969 documenting the risk of maternal alcohol consumption among the offspring of 11 alcoholic mothers. This study is arguably the true source of the modern understanding. The infants were studied by dysmorphologists Kenneth Lyons Jones and David Weyhe Smith, colleagues of Ulleland at University of Washington, who identified a pattern of "craniofacial, limb, and cardiovascular defects associated with prenatal onset growth deficiency and developmental delay" in eight children. The pattern of malformations indicated that the damage was prenatal. They named these defects "fetal alcohol syndrome". News of the discovery shocked some, while others were skeptical of the findings. While many syndromes are eponymous, i.e. named after the physician first reporting the association of symptoms, Smith named FAS after the causal agent of the symptoms. He reasoned that doing so would encourage prevention, believing that if people knew maternal alcohol consumption caused the syndrome, then abstinence during pregnancy would follow from patient education and public awareness. At the time, nobody was aware of the full range of possible birth defects from FAS or its rate of prevalence.
In 1978, within nine years of the Washington discovery, animal studies, including non-human monkey studies carried out at the University of Washington Primate Center by Sterling Clarren, had confirmed that alcohol was a teratogen. By 1978, 245 cases of FAS had been reported by medical researchers, and the syndrome began to be described as the most frequent known cause of intellectual disability. In 1979, the Washington and Nantes findings were confirmed by a research group in Gothenburg, Sweden. Researchers in France, Sweden, and the United States were struck by how similar these children looked, though they were not related, and how they behaved in the same unfocused and hyperactive manner.
"Spectrum" rather than Syndrome
Over time, subsequent research and clinical experience suggested that a range of effects could arise from prenatal alcohol exposure. The term fetal alcohol effects (FAE) was used for alcohol-related neurodevelopmental disorder and alcohol-related birth defects. It was initially used in research studies to describe humans and animals in whom teratogenic effects were seen after confirmed prenatal alcohol exposure (or unknown exposure for humans), but without obvious physical anomalies. Smith (1981) described FAE as an "extremely important concept" to highlight the debilitating effects of brain damage, regardless of the growth or facial features. This term fell out of favor with clinicians in the 1990s because it was often regarded by the public as a less severe disability than FAS, when in fact its effects could be just as detrimental. In 1996 the replacement terms ARBD and ARND were introduced. In 2002, the US Congress mandated that the CDC develop diagnostic guidelines for FAS and in 2004 a definition of a term that already had been used by some in the nineties, the Fetal Alcohol Spectrum Disorder (FASD) was adopted,> to include FAS as well as other conditions resulting from prenatal alcohol exposure.
Currently, FAS is the only expression of prenatal alcohol exposure defined by the International Statistical Classification of Diseases and Related Health Problems and assigned ICD-9 and diagnoses.
Alcohol-related birth defects (ARBD), formerly known as possible fetal alcohol effect (PFAE), was a term proposed as an alternative to FAE and PFAE. The IOM presents ARBD as a list of congenital anomalies that are linked to maternal alcohol use but have no key features of FASD. PFAE and ARBD have fallen out of favor because these anomalies are not necessarily specific to maternal alcohol consumption and are not criteria for diagnosis of FASD.
In 2013, the American Psychiatric Association introduced neurobehavioral disorder associated with prenatal alcohol exposure (ND-PAE).
Society and culture
Criminalization
Criminalization of substance use during pregnancy because of harm to the fetus or child is fiercely debated. Elizabeth Armstrong has questioned the zero-tolerance approach taken towards alcohol consumption during pregnancy, describing it as a moral panic. While heavy alcohol consumption during pregnancy is known to be damaging to the unborn child, the effects of low intakes remain debatable, particularly in the absence of randomized controlled trials (c.f. ). The UK's abstinence recommendation was not chosen based on scientific evidence, but rather because it was simple advice that would ensure no one underestimated the risk. Tennessee's 2014 fetal assault law (which expired in 2016) was criticized for not addressing alcohol use. The law criminalized opioid use during pregnancy and resulted in women avoiding professional medical care for fear of prosecution. A wide variety of professional organizations oppose criminalization. Minnesota, North Dakota, Oklahoma, South Dakota, and Wisconsin have statutory authorization for the involuntary civil commitment of women who abuse alcohol during pregnancy.
2016 CDC controversy
In 2016, a CDC press release and infographic entitled "More than 3 million US women at risk for alcohol-exposed pregnancy" caused controversy. The CDC release contained the message "The risk is real. Why take the chance?". Darlena Cunha of Times Magazine interpreted the infographic as telling all women of child-bearing age not to drink at all, in case they might accidentally fall pregnant, and called them "scare tactics" and "shaming recommendations". Julie Beck said that the infographic insinuated that "your womb is a Schrodinger's box and you shouldn't pour alcohol into it unless you've peeked in there to be 100 percent sure the coast is clear". The CDC later clarified that the infographic was not intended to make any new guidelines or recommendations for women who are pre-pregnant, but rather to encourage conversations about alcohol with health professionals. Nonetheless, half of the pregnancies in developed countries and over 80% in developing countries are unplanned. Many women do not realize they are pregnant during the early stages and continue drinking when pregnant.
In fiction
In Aldous Huxley's 1932 novel Brave New World (where all fetuses are gestated in vitro in a factory), lower caste fetuses are created by receiving alcohol transfusions (Bokanovsky Process) to reduce intelligence and height, thus conditioning them for simple, menial tasks. Connections between alcohol and incubating embryos are made multiple times in the novel.
The main character of the 2009 film Defendor is implied to have the condition.
Tony Loneman, a character in Tommy Orange's 2018 novel There There, was born with fetal alcohol syndrome, which he calls "the Drome".
| Biology and health sciences | Disabilities | Health |
1249122 | https://en.wikipedia.org/wiki/Water%20supply%20network | Water supply network | A water supply network or water supply system is a system of engineered hydrologic and hydraulic components that provide water supply. A water supply system typically includes the following:
A drainage basin (see water purification – sources of drinking water)
A raw water collection point (above or below ground) where the water accumulates, such as a lake, a river, or groundwater from an underground aquifer. Raw water may be transferred using uncovered ground-level aqueducts, covered tunnels, or underground water pipes to water purification facilities.
Water purification facilities. Treated water is transferred using water pipes (usually underground).
Water storage facilities such as reservoirs, water tanks, or water towers. Smaller water systems may store the water in cisterns or pressure vessels. Tall buildings may also need to store water locally in pressure vessels in order for the water to reach the upper floors.
Additional water pressurizing components such as pumping stations may need to be situated at the outlet of underground or aboveground reservoirs or cisterns (if gravity flow is impractical).
A pipe network for distribution of water to consumers (which may be private houses or industrial, commercial, or institution establishments) and other usage points (such as fire hydrants)
Connections to the sewers (underground pipes, or aboveground ditches in some developing countries) are generally found downstream of the water consumers, but the sewer system is considered to be a separate system, rather than part of the water supply system.
Water supply networks are often run by public utilities of the water industry.
Water abstraction and raw water transfer
Raw water (untreated) is from a surface water source (such as an intake on a lake or a river) or from a groundwater source (such as a water well drawing from an underground aquifer) within the watershed that provides the water resource.
The raw water is transferred to the water purification facilities using uncovered aqueducts, covered tunnels or underground water pipes.
Water treatment
Virtually all large systems must treat the water; a fact that is tightly regulated by global, state and federal agencies, such as the World Health Organization (WHO) or the United States Environmental Protection Agency (EPA). Water treatment must occur before the product reaches the consumer and afterwards (when it is discharged again). Water purification usually occurs close to the final delivery points to reduce pumping costs and the chances of the water becoming contaminated after treatment.
Traditional surface water treatment plants generally consists of three steps: clarification, filtration and disinfection. Clarification refers to the separation of particles (dirt, organic matter, etc.) from the water stream. Chemical addition (i.e. alum, ferric chloride) destabilizes the particle charges and prepares them for clarification either by settling or floating out of the water stream. Sand, anthracite or activated carbon filters refine the water stream, removing smaller particulate matter. While other methods of disinfection exist, the preferred method is via chlorine addition. Chlorine effectively kills bacteria and most viruses and maintains a residual to protect the water supply through the supply network.
Water distribution network
The product, delivered to the point of consumption, is called potable water if it meets the water quality standards required for human consumption.
The water in the supply network is maintained at positive pressure to ensure that water reaches all parts of the network, that a sufficient flow is available at every take-off point and to ensure that untreated water in the ground cannot enter the network. The water is typically pressurised by pumping the water into storage tanks constructed at the highest local point in the network. One network may have several such service reservoirs.
In small domestic systems, the water may be pressurised by a pressure vessel or even by an underground cistern (the latter however does need additional pressurizing). This eliminates the need of a water tower or any other heightened water reserve to supply the water pressure.
These systems are usually owned and maintained by local governments such as cities or other public entities, but are occasionally operated by a commercial enterprise (see water privatization). Water supply networks are part of the master planning of communities, counties, and municipalities. Their planning and design requires the expertise of city planners and civil engineers, who must consider many factors, such as location, current demand, future growth, leakage, pressure, pipe size, pressure loss, fire fighting flows, etc.—using pipe network analysis and other tools.
As water passes through the distribution system, the water quality can degrade by chemical reactions and biological processes. Corrosion of metal pipe materials in the distribution system can cause the release of metals into the water with undesirable aesthetic and health effects. Release of iron from unlined iron pipes can result in customer reports of "red water" at the tap. Release of copper from copper pipes can result in customer reports of "blue water" and/or a metallic taste. Release of lead can occur from the solder used to join copper pipe together or from brass fixtures. Copper and lead levels at the consumer's tap are regulated to protect consumer health.
Utilities will often adjust the chemistry of the water before distribution to minimize its corrosiveness. The simplest adjustment involves control of pH and alkalinity to produce a water that tends to passivate corrosion by depositing a layer of calcium carbonate. Corrosion inhibitors are often added to reduce release of metals into the water. Common corrosion inhibitors added to the water are phosphates and silicates.
Maintenance of a biologically safe drinking water is another goal in water distribution. Typically, a chlorine based disinfectant, such as sodium hypochlorite or monochloramine is added to the water as it leaves the treatment plant. Booster stations can be placed within the distribution system to ensure that all areas of the distribution system have adequate sustained levels of disinfection.
Topologies
Like electric power lines, roads, and microwave radio networks, water systems may have a loop or branch network topology, or a combination of both. The piping networks are circular or rectangular. If any one section of water distribution main fails or needs repair, that section can be isolated without disrupting all users on the network.
Most systems are divided into zones. Factors determining the extent or size of a zone can include hydraulics, telemetry systems, history, and population density. Sometimes systems are designed for a specific area then are modified to accommodate development. Terrain affects hydraulics and some forms of telemetry. While each zone may operate as a stand-alone system, there is usually some arrangement to interconnect zones in order to manage equipment failures or system failures.
Water network maintenance
Water supply networks usually represent the majority of assets of a water utility. Systematic documentation of maintenance works using a computerized maintenance management system (CMMS) is a key to a successful operation of a water utility.
Sustainable urban water supply
A sustainable urban water supply network covers all the activities related to provision of potable water. Sustainable development is of increasing importance for the water supply to urban areas. Incorporating innovative water technologies into water supply systems improves water supply from sustainable perspectives. The development of innovative water technologies provides flexibility to the water supply system, generating a fundamental and effective means of sustainability based on an integrated real options approach.
Water is an essential natural resource for human existence. It is needed in every industrial and natural process, for example, it is used for oil refining, for liquid-liquid extraction in hydro-metallurgical processes, for cooling, for scrubbing in the iron and the steel industry, and for several operations in food processing facilities.
It is necessary to adopt a new approach to design urban water supply networks; water shortages are expected in the forthcoming decades and environmental regulations for water utilization and waste-water disposal are increasingly stringent.
To achieve a sustainable water supply network, new sources of water are needed to be developed, and to reduce environmental pollution.
The price of water is increasing, so less water must be wasted and actions must be taken to prevent pipeline leakage. Shutting down the supply service to fix leaks is less and less tolerated by consumers. A sustainable water supply network must monitor the freshwater consumption rate and the waste-water generation rate.
Many of the urban water supply networks in developing countries face problems related to population increase, water scarcity, and environmental pollution.
Population growth
In 1900 just 13% of the global population lived in cities. By 2005, 49% of the global population lived in urban areas. In 2030 it is predicted that this statistic will rise to 60%. Attempts to expand water supply by governments are costly and often not sufficient. The building of new illegal settlements makes it hard to map, and make connections to, the water supply, and leads to inadequate water management. In 2002, there were 158 million people with inadequate water supply. An increasing number of people live in slums, in inadequate sanitary conditions, and are therefore at risk of disease.
Water scarcity
Potable water is not well distributed in the world. 1.8 million deaths are attributed to unsafe water supplies every year, according to the WHO. Many people do not have any access, or do not have access to quality and quantity of potable water, though water itself is abundant. Poor people in developing countries can be close to major rivers, or be in high rainfall areas, yet not have access to potable water at all. There are also people living where lack of water creates millions of deaths every year.
Where the water supply system cannot reach the slums, people manage to use hand pumps, to reach the pit wells, rivers, canals, swamps and any other source of water. In most cases the water quality is unfit for human consumption. The principal cause of water scarcity is the growth in demand. Water is taken from remote areas to satisfy the needs of urban areas. Another reason for water scarcity is climate change: precipitation patterns have changed; rivers have decreased their flow; lakes are drying up; and aquifers are being emptied.
Governmental issues
In developing countries many governments are corrupt and poor and they respond to these problems with frequently changing policies and non clear agreements. Water demand exceeds supply, and household and industrial water supplies are prioritised over other uses, which leads to water stress. Potable water has a price in the market; water often becomes a business for private companies, which earn a profit by putting a higher price on water, which imposes a barrier for lower-income people. The Millennium Development Goals propose the changes required.
Goal 6 of the United Nations' Sustainable Development Goals is to "Ensure availability and sustainable management of water and sanitation for all". This is in recognition of the human right to water and sanitation, which was formally acknowledged at the United Nations General Assembly in 2010, that "clean drinking water and sanitation are essential to the recognition of all human rights". Sustainable water supply includes ensuring availability, accessibility, affordability and quality of water for all individuals.
In advanced economies, the problems are about optimising existing supply networks. These economies have usually had continuing evolution, which allowed them to construct infrastructure to supply water to people. The European Union has developed a set of rules and policies to overcome expected future problems.
There are many international documents with interesting, but not very specific, ideas and therefore they are not put into practice. Recommendations have been made by the United Nations, such as the Dublin Statement on Water and Sustainable Development.
Optimizing the water supply network
The yield of a system can be measured by either its value or its net benefit. For a water supply system, the true value or the net benefit is a reliable water supply service having adequate quantity and good quality of the product. For example, if the existing water supply of a city needs to be extended to supply a new municipality, the impact of the new branch of the system must be designed to supply the new needs, while maintaining supply to the old system.
Single-objective optimization
The design of a system is governed by multiple criteria, one being cost. If the benefit is fixed, the least cost design results in maximum benefit. However, the least cost approach normally results in a minimum capacity for a water supply network. A minimum cost model usually searches for the least cost solution (in pipe sizes), while satisfying the hydraulic constraints such as: required output pressures, maximum pipe flow rate and pipe flow velocities. The cost is a function of pipe diameters; therefore the optimization problem consists of finding a minimum cost solution by optimising pipe sizes to provide the minimum acceptable capacity.
Multi-objective optimization
However, according to the authors of the paper entitled, “Method for optimizing design and rehabilitation of water distribution systems”, “the least capacity is not a desirable solution to a sustainable water supply network in a long term, due to the uncertainty of the future demand”. It is preferable to provide extra pipe capacity to cope with unexpected demand growth and with water outages. The problem changes from a single objective optimization problem (minimizing cost), to a multi-objective optimization problem (minimizing cost and maximizing flow capacity).
Weighted sum method
To solve a multi-objective optimization problem, it is necessary to convert the problem into a single objective optimization problem, by using adjustments, such as a weighted sum of objectives, or an ε-constraint method. The weighted sum approach gives a certain weight to the different objectives, and then factors in all these weights to form a single objective function that can be solved by single factor optimization. This method is not entirely satisfactory, because the weights cannot be correctly chosen, so this approach cannot find the optimal solution for all the original objectives.
The constraint method
The second approach (the constraint method), chooses one of the objective functions as the single objective, and the other objective functions are treated as constraints with a limited value. However, the optimal solution depends on the pre-defined constraint limits.
Sensitivity analysis
The multiple objective optimization problems involve computing the tradeoff between the costs and benefits resulting in a set of solutions that can be used for sensitivity analysis and tested in different scenarios. But there is no single optimal solution that will satisfy the global optimality of both objectives. As both objectives are to some extent contradictory, it is not possible to improve one objective without sacrificing the other. It is necessary in some cases use a different approach. (e.g. Pareto Analysis), and choose the best combination.
Operational constraints
Returning to the cost objective function, it cannot violate any of the operational constraints. Generally this cost is dominated by the energy cost for pumping. “The operational constraints include the standards of customer service, such as: the minimum delivered pressure, in addition to the physical constraints such as the maximum and the minimum water levels in storage tanks to prevent overtopping and emptying respectively.”
In order to optimize the operational performance of the water supply network, at the same time as minimizing the energy costs, it is necessary to predict the consequences of different pump and valve settings on the behavior of the network.
Apart from Linear and Non-linear Programming, there are other methods and approaches to design, to manage and operate a water supply network to achieve sustainability—for instance, the adoption of appropriate technology coupled with effective strategies for operation and maintenance. These strategies must include effective management models, technical support to the householders and industries, sustainable financing mechanisms, and development of reliable supply chains. All these measures must ensure the following: system working lifespan; maintenance cycle; continuity of functioning; down time for repairs; water yield and water quality.
Sustainable development
In an unsustainable system there is insufficient maintenance of the water networks, especially in the major pipe lines in urban areas. The system deteriorates and then needs rehabilitation or renewal.
Householders and sewage treatment plants can both make the water supply networks more efficient and sustainable. Major improvements in eco-efficiency are gained through systematic separation of rainfall and wastewater. Membrane technology can be used for recycling wastewater.
The municipal government can develop a “Municipal Water Reuse System” which is a current approach to manage the rainwater. It applies a water reuse scheme for treated wastewater, on a municipal scale, to provide non-potable water for industry, household and municipal uses. This technology consists in separating the urine fraction of sanitary wastewater, and collecting it for recycling its nutrients. The feces and graywater fraction is collected, together with organic wastes from the households, using a gravity sewer system, continuously flushed with non-potable water. The water is treated anaerobically and the biogas is used for energy production.
One effective way to achieve sustainable water management is to shift emphasis towards decentralized water projects, such as drip irrigation diffusion in India. This project covers large spatial areas while relying on individual technological adoption decisions, offering scalable solutions that can mitigate water scarcity and enhance agricultural productivity.
Another method that can be utilized is through the promoting of community engagement and resistance against unsustainable water infrastructure projects. Grassroots movements, as observed in anti-dam protests in various countries, play a crucial role in challenging dominant development narratives and advocating for more socially and ecologically just water management practices.
Municipalities and other forms of local governments should also invest in innovative technologies, such as membrane technology for wastewater recycling, and develop policy frameworks that incentivize eco-efficient practices. Municipal water reuse systems, as demonstrated in implementation, offer promising avenues for integrating wastewater treatment and resource recovery into urban water networks.
The sustainable water supply system is an integrated system including water intake, water utilization, wastewater discharge and treatment and water environmental protection. It requires reducing freshwater and groundwater usage in all sectors of consumption. Developing sustainable water supply systems is a growing trend, because it serves people's long-term interests. There are several ways to reuse and recycle the water, in order to achieve long-term sustainability, such as:
Gray water re-use and treatment: gray water is wastewater coming from baths, showers, sinks and washbasins. If this water is treated it can be used as a source of water for uses other than drinking. Depending on the type of gray water and its level of treatment, it can be re-used for irrigation and toilet flushing. According to an investigation about the impacts of domestic grey water reuse on public health, carried out by the New South Wales Health Centre in Australia in the year 2000, grey water contains less nitrogen and fecal pathogenic organisms than sewage, and the organic content of grey water decomposes more rapidly.
Ecological treatment systems use little energy: there are many applications in gray water re-use, such as reed beds, soil treatment systems and plant filters. This process is ideal for gray water re-use, because of easier maintenance and higher removal rates of organic matter, ammonia, nitrogen and phosphorus.
Other possible approaches to scoping models for water supply, applicable to any urban area, include the following:
Sustainable drainage system
Borehole extraction
Intercluster groundwater flow
Canal and river extraction
Aquifer storage
A more user-friendly indoor water use
The Dublin Statement on Water and Sustainable Development is a good example of the new trend to overcome water supply problems. This statement, suggested by advanced economies, has come up with some principles that are of great significance to urban water supply. These are:
Fresh water is a finite and vulnerable resource, essential to sustain life, development and the environment.
Water development and management should be based on a participatory approach, involving users, planners and policy-makers at all levels.
Women play a central part in the provision, management and safeguarding of water. Institutional arrangements should reflect the role of women in water provision and protection.
Water has an economic value in all its competing uses and should be recognized as an economic good.
From these statements, developed in 1992, several policies have been created to give importance to water and to move urban water system management towards sustainable development. The Water Framework Directive by the European Commission is a good example of what has been created there out of former policies.
Future approaches
There is great need for a more sustainable water supply systems. To achieve sustainability several factors must be tackled at the same time: climate change, rising energy cost, and rising populations. All of these factors provoke change and put pressure on management of available water resources.
An obstacle to transforming conventional water supply systems, is the amount of time needed to achieve the transformation. More specifically, transformation must be implemented by municipal legislation bodies, which always need short-term solutions too. Another obstacle to achieving sustainability in water supply systems is the insufficient practical experience with the technologies required, and the missing know-how about the organization and the transition process.
Urban water infrastructure faces several challenges that undermine its sustainability and resilience. One critical issue highlighted in recent research is the vulnerability of water networks to climate variability and extreme weather events. Poor seasonal rains, as observed in the case of the Panama Canal's lock and dam infrastructure, exemplify how inadequate water supply can strain water-intensive infrastructure, raising questions about engineering legitimacy and the reliability of water systems.
Another key challenge is the unequal development associated with large-scale water infrastructure projects such as dams and canals . Such projects, while aimed at promoting economic growth, often actually reproduce social and economic inequalities by displacing rural communities and marginalizing indigenous populations. This phenomenon of "accumulation by dispossession" further emphasizes the need for more equitable and inclusive approaches to water infrastructure development.
Possible ways to improve this situation is simulating of the network, implementing pilot projects, learning from the costs involved and the benefits achieved.
| Technology | Food, water and health | null |
15690011 | https://en.wikipedia.org/wiki/Staphylococcal%20infection | Staphylococcal infection | A staphylococcal infection or staph infection is an infection caused by members of the Staphylococcus genus of bacteria.
These bacteria commonly inhabit the skin and nose where they are innocuous, but may enter the body through cuts or abrasions which may be nearly invisible. Once inside the body, the bacteria may spread to a number of body systems and organs, including the heart, where the toxins produced by the bacteria may cause cardiac arrest. Once the bacterium has been identified as the cause of the illness, treatment is often in the form of antibiotics and, where possible, drainage of the infected area. However, many strains of this bacterium have become antibiotic resistant; for those with these kinds of infection, the body's own immune system is the only defense against the disease. If that system is weakened or compromised, the disease may progress rapidly. Anyone can contract staph, but pregnant women, children, and people with chronic diseases or who are immuno-deficient are often more susceptible to contracting an infection.
Types
Other infections include:
Closed-space infections of the fingertips, known as paronychia.
Suspected involvement in atopic dermatitis (eczema), including related clinical trials.
Coagulase-positive
The main coagulase-positive staphylococcus is Staphylococcus aureus, although not all strains of Staphylococcus aureus are coagulase positive. These bacteria can survive on dry surfaces, increasing the chance of transmission. S. aureus is also implicated in toxic shock syndrome; during the 1980s some tampons allowed the rapid growth of S. aureus, which released toxins that were absorbed into the bloodstream. Any S. aureus infection can cause the staphylococcal scalded skin syndrome, a cutaneous reaction to exotoxin absorbed into the bloodstream. It can also cause a type of septicaemia called pyaemia. The infection can be life-threatening. Problematically, methicillin-resistant Staphylococcus aureus (MRSA) has become a major cause of hospital-acquired infections. MRSA has also been recognized with increasing frequency in community-acquired infections. The symptoms of a staphylococcal infection include a collection of pus, such as a boil or furuncle, or abscess. The area is typically tender or painful and may be reddened or swollen.
Coagulase-negative
S. epidermidis, a coagulase-negative staphylococcus species, is a commensal of the skin, but can cause severe infections in immune-suppressed patients and those with central venous catheters.
S. saprophyticus, another coagulase-negative species that is part of the normal vaginal flora, is predominantly implicated in uncomplicated lower genitourinary tract infections in young sexually active women.
Other staphylococcal species have been implicated in human infections, notably S. lugdunensis, S. schleiferi, and S. caprae.
Causes
Staph infections have a multitude of different causes, such as:
Open wounds – This is by far the biggest cause of staph infection. Any open wound, even ones as small as a paper cut, are vulnerable to being infected. Staph bacteria will enter the body through any open wound, so it is important to properly treat, disinfect, and bandage any wounds.
Contact with infected persons or surfaces – Staph infections are very contagious when in contact with a person that is already infected. A person with staph infection is contagious until the bacteria are completely out of their body, and any wounds from the infection are healed. It is common to see the spread of staph in contact sports; i.e. wrestling, through contact in locker rooms, or by sharing any equipment.
Weakened immune system – Anyone with a weakened immune system for any reason can be more easily affected by staph bacteria, because their bodies are unable to defend against infectious bacteria as well.
Unwashed linens – Staph bacteria are very resistant under harsh conditions, and they will cling to objects where they can create a niche. Unwashed bath towels, blanket, bed sheets, and clothes can create a great environment for these bacteria to grow. This is important to recognize, because every single day people use linens in their daily lives.
Infection after surgery – Hospitals are a very common place for staph bacteria to contaminate. This becomes problematic when people are in surgery, because in some cases staph can be introduced to a person's body when an incision is opened.
Invasive devices – Medical devices that have any connection to organs to the outside of the body are very problematic, because they allow an easy open pathway into the body. Examples of these devices are; catheters, dialysis tubing, feeding tubes, breathing tubes, etc.
Signs and symptoms
Staph infection is typically characterized by redness, pus, swelling, and tenderness in areas of the infection. But, each type of skin infection caused by staph bacteria is different.
A few common skin infections caused by staph bacteria are:
Boils – Boils are the most common type of staph infection, they are pockets of white pus that start where a hair follicle or oil gland is. The boil is tender and red where the infection is located on the skin.
Impetigo – Impetigo is most prominent among children, and is usually located around their mouth, nose, hands, and feet. It shows up like a rash of painful blisters, will eventually produce pus that is yellowish in color.
Cellulitis – Cellulitis is also rash-like; the skin that is infected will be red, swollen, and usually warm to the touch. Cellulitis commonly infects the lower legs, but can also, less commonly, affect the face and arms.
Staphylococcus scalded skin syndrome – Staphylococcus scalded skin syndrome is caused by toxins produced when a staph infection gets too severe. It is characterized by a fever, rash, and blisters.
Methicillin-resistant Staphylococcus aureus (MRSA) – MRSA is one of the most common antibiotic-resistant strains of staph bacteria. It is more difficult to treat than other staph infections. MRSA causes rashes, boils, sores, and other abscesses.
Bacterial identification
In the microbiology lab, Staphylococcus is mainly suspected when seeing Gram-positive cocci in clusters.
Treatment
Treatment for staph infection varies depending on the type and severity of infection. Common treatments are antibiotics, topical creams, and drainage/cleaning of infectious wounds.
Etymology
The generic name Staphylococcus is derived from the Greek word "staphyle", meaning bunch of grapes, and "kokkos", meaning granule. The bacteria, when seen under a microscope, appear like a branch of grapes or nuts.
Epidemiology
Staphylococcus bacteria is one of the leading community-acquired bacteria. According to the CDC, after a push from hospitals to better prevent staph infections, the percentage of people affected has dropped dramatically. However, staph infections are still prominent and a cause for concern among healthcare professionals, especially new antibiotic-resistant strains. In the U.S., the incidence of staph infection is around 38.2 to 45.7 per 100,000 person-years, whereas other First World countries have an average incidence rate of 10 to 30 per 100,000 person-years.
| Biology and health sciences | Bacterial infections | Health |
22293372 | https://en.wikipedia.org/wiki/SAR%20supergroup | SAR supergroup | SAR or Harosa is a highly diverse clade of eukaryotes, often considered a supergroup, that includes stramenopiles (heterokonts), alveolates, and rhizarians. It is a node-based taxon, including all descendants of the three groups' last common ancestor, and comprises most of the now-rejected Chromalveolata. Their sister group has been found to be telonemids, with which they make up the TSAR clade.
Etymology
The name SAR is an acronym derived from the first letters of its three constituent clades; it has been alternatively spelled "RAS". The term "Harosa" (at the subkingdom level) has also been used, with Stramenopiles replaced by its synonym Heterokonta in this variant of the acronym.
History of discovery
Before the discovery of the SAR supergroup, stramenopiles and alveolates were classified in the supergroup Chromalveolata alongside haptophytes and cryptomonads, being believed to have acquired plastids through secondary endosymbiosis of red algae through a common ancestor. Meanwhile, Rhizaria was traditionally considered to be a separate supergroup. More recent phylogenetic studies confirmed that stramenopiles and alveolates diverged with rhizarians as part of the SAR lineage. This clade has been found by later phylogenomic studies to be robustly characterized compared to other supergroups.
This groups excludes haptophytes and cryptomonads, hypothesized to have acquired plastids in separate endosymbiosis events, leading Okamoto et al. (2009) to propose the clade Hacrobia to accommodate them.
Diversity
The SAR supergroup encompasses a variety of morphologies and ecological niches, from microscopic zoo- and phytoplankton to massive kelp forests. The group includes both photosynthetic and non-photosynthetic forms. Photosynthesis arose independently across various stramenopiles and alveolates lineages through secondary or higher-order endosymbiosis events, acquiring plastids of red algal origin, while chlorarachniophyte rhizarians captured plastids from green algae, retaining vestigial nucleomorphs.
It has been estimated that SAR encompasses up to half of all eukaryotic diversity.
Owing to the clade's discovery through phylogenomics, there are no known synapomorphies uniting its various members. This was already the case for its subclade Rhizaria, established earlier through similar means. On the other hand, Stramenopiles is well-defined morphologically, characterized by an anterior flagellum with tripartite bristles (mastigonemes), while Alveolata is united by the presence of cortical alveoli.
Nonetheless, studies of telonemids, believed to be the sister group to SAR, have revealed characteristics such as tripartite hair and peripheral vacuoles, potentially homologous to similar structures in stramenopiles and alveolates. This brings into light the possibility of these structures being ancestrally shared by the clade, with cortical alveoli originating from peripheral vacuoles under this hypothesis.
Internal phylogeny
A 2021 analysis places Alveolata and Stramenopiles in Halvaria, as sister to Rhizaria.
| Biology and health sciences | SAR supergroup | Plants |
22295187 | https://en.wikipedia.org/wiki/Epoch%20%28computing%29 | Epoch (computing) | In computing, an epoch is a fixed date and time used as a reference from which a computer measures system time. Most computer systems determine time as a number representing the seconds removed from a particular arbitrary date and time. For instance, Unix and POSIX measure time as the number of seconds that have passed since Thursday 1 January 1970 00:00:00 UT, a point in time known as the Unix epoch. The C# programming language and Windows NT systems up to and including Windows 11 and Windows Server 2022 measure time as the number of 100-nanosecond intervals that have passed since 00:00:00 UTC on 1 January in the years AD 1 and AD 1601, respectively, making those points in time the epochs for those systems.
Computing epochs are almost always specified as midnight Universal Time on some particular date.
Resolution and representation
Software timekeeping systems vary widely in the resolution of time measurement; some systems may use time units as large as a day, while others may use nanoseconds. For example, for an epoch date of midnight UTC (00:00) on 1 January 1900, and a time unit of a second, the time of the midnight (24:00) between 1 January 1900 and 2 January 1900 is represented by the number 86400, the number of seconds in one day. When times prior to the epoch need to be represented, it is common to use the same system, but with negative numbers.
Such representation of time is mainly for internal use. On systems where date and time are important in the human sense, software will almost always convert this internal number into a date and time representing a human calendar.
Problems
Computers do not generally store arbitrarily large numbers. Instead, each number stored by a computer is allotted a fixed amount of space. Therefore, when the number of time units that have elapsed since a system's epoch exceeds the largest number that can fit in the space allotted to the time representation, the time representation overflows, and problems can occur. While a system's behavior after overflow occurs is not necessarily predictable, in most systems the number representing the time will reset to zero, and the computer system will think that the current time is the epoch time again.
Most famously, older systems that counted time as the number of years elapsed since the epoch of 1 January 1900 and which only allotted enough space to store the numbers 0 through 99, experienced the Year 2000 problem. These systems (if not corrected beforehand) would interpret the date 1 January 2000 as 1 January 1900, leading to unpredictable errors at the beginning of the year 2000.
Even systems that allocate more storage to the time representation are not immune from this kind of error. Many Unix-like operating systems which keep time as seconds elapsed from the epoch date of 1 January 1970, and allot timekeeping enough storage to store numbers as large as will experience an overflow problem on 19 January 2038. This is known as the Year 2038 problem.
Other more subtle timekeeping problems exist in computing, such as accounting for leap seconds, which are not observed with any predictability or regularity. Additionally, applications that need to represent historical dates and times (for example, representing a date prior to the switch from the Julian calendar to the Gregorian calendar) must use specialized timekeeping libraries.
Finally, some software must maintain compatibility with older software that does not keep time in strict accordance with traditional timekeeping systems. For example, Microsoft Excel observes the fictional date of 29 February 1900 in order to maintain bug compatibility with older versions of Lotus 1-2-3. Lotus 1-2-3 observed the date due to an error; by the time the error was discovered, it was too late to fix it—"a change now would disrupt formulas which were written to accommodate this anomaly".
In satellite-based time systems
There are at least six satellite navigation systems, all of which function by transmitting time signals. Of the only two satellite systems with global coverage, GPS calculates its time signal from an epoch, whereas GLONASS calculates time as an offset from UTC, with the UTC input adjusted for leap seconds. Of the only two other systems aiming for global coverage, Galileo calculates from an epoch and BeiDou calculates from UTC without adjustment for leap seconds. GPS also transmits the offset between UTC time and GPS time and must update this offset every time there is a leap second, requiring GPS-receiving devices to handle the update correctly. In contrast, leap seconds are transparent to GLONASS users.
The complexities of calculating UTC from an epoch are explained by the European Space Agency in Galileo documentation under "Equations to correct system timescale to reference timescale".
Notable epoch dates in computing
The following table lists epoch dates used by popular software and other computer-related systems. The time in these systems is stored as the quantity of a particular time unit (days, seconds, nanoseconds, etc.) that has elapsed since a stated time (usually midnight UTC at the beginning of the given date).
| Technology | Computer architecture concepts | null |
1860870 | https://en.wikipedia.org/wiki/Environmental%20justice | Environmental justice | Environmental justice is a social movement that addresses injustice that occurs when poor or marginalized communities are harmed by hazardous waste, resource extraction, and other land uses from which they do not benefit. The movement has generated hundreds of studies showing that exposure to environmental harm is inequitably distributed.
The movement began in the United States in the 1980s. It was heavily influenced by the American civil rights movement and focused on environmental racism within rich countries. The movement was later expanded to consider gender, international environmental injustice, and inequalities within marginalized groups. As the movement achieved some success in rich countries, environmental burdens were shifted to the Global South (as for example through extractivism or the global waste trade). The movement for environmental justice has thus become more global, with some of its aims now being articulated by the United Nations. The movement overlaps with movements for Indigenous land rights and for the human right to a healthy environment.
The goal of the environmental justice movement is to achieve agency for marginalized communities in making environmental decisions that affect their lives. The global environmental justice movement arises from local environmental conflicts in which environmental defenders frequently confront multi-national corporations in resource extraction or other industries. Local outcomes of these conflicts are increasingly influenced by trans-national environmental justice networks.
Environmental justice scholars have produced a large interdisciplinary body of social science literature that includes contributions to political ecology, environmental law, and theories on justice and sustainability.
Definitions
The United States Environmental Protection Agency defines environmental justice as:the fair treatment and meaningful involvement of all people regardless of race, color, national origin, or income with respect to the development, implementation, and enforcement of environmental laws, regulations and policies. Fair treatment means that no group of people, including racial, ethnic, or socio-economic groups, should bear a disproportionate share of the negative environmental consequences resulting from industrial, municipal, and commercial operations or the execution of federal, state, local, and tribal programs and policies Environmental justice is also discussed as environmental racism or environmental inequality.
Environmental justice is typically defined as distributive justice, which is the equitable distribution of environmental risks and benefits. Some definitions address procedural justice, which is the fair and meaningful participation in decision-making. Other scholars emphasise recognition justice, which is the recognition of oppression and difference in environmental justice communities. People's capacity to convert social goods into a flourishing community is a further criteria for a just society. However, initiatives have been taken to expand the notion of environmental justice beyond the three pillars of distribution, participation, and recognition to also include the dimensions of self-governing authority, relational ontologies, and epistemic justice.
Robert D. Bullard writes that environmental justice, as a social movement and ideological stewardship, may instead be seen as a conversation of equity. Bullard writes that equity is distilled into three board categories: procedural, geographic, and social. From his publication “Confronting Environmental Racism in the Twenty-First Century,” he draws our the difference between the three within the context of environmental injustices:Procedural equity refers to the “fairness” question: the extent that rules, regulations, evaluation criteria and enforcement are applied uniformly across the board and in a non-discriminatory way. Unequal protection might result from nonscientific and undemocratic decisions, exclusionary practices, public hearings held in remote locations and at inconvenient times, and use of English-only material as the language in which to communicate and conduct hearings for non-English-speaking publics.
Geographic equity refers to the location and spatial configuration of communities and their proximity to environmental hazards, noxious facilities and locally unwanted land uses (Lulus) such as landfills, incinerators, sewage treatment plants, lead smelters, refineries and other noxious facilities. For example, unequal protection may result from land-use decisions that determine the location of residential amenities and disamenities. The poor and communities of colour often suffer a “triple” vulnerability of noxious facility siting, as do the unincorporated—sparsely populated communities that are not legally chartered as cities or municipalities and are therefore usually governed by distant county governments rather than having their own locally elected officials.
Social equity assesses the role of sociological factors (race, ethnicity, class, culture, life styles, political power, etc.) on environmental decision making. Poor people and people of colour often work in the most dangerous jobs and live in the most polluted neighbourhoods, their children exposed to all kinds of environmental toxins in the playgrounds and in their homes.
Indigenous environmental justice
In non-Native communities, where toxic industries and other discriminatory practices are disproportionately occurring, residents rely on laws and statutory frameworks outlined by the EPA. They rely on distributive justice, centered around the nature of private property. Native Americans do not fall under the same statutory frameworks as they are citizens of Indigenous nations, not ethnic minorities. As individuals, they are subject to American laws. As nations, they are subject to a separate legal regime, constructed on the basis of pre-existing sovereignty acknowledged by treaty and the U.S. Constitution. Environmental justice to Indigenous persons is not understood by legal entities but rather their distinct cultural and religious doctrines.
Environmental Justice for Indigenous peoples follows a model that frames issues in terms of their colonial condition and can affirm decolonization as a potential framework within environmental justice. While Indigenous peoples’ lived experiences vary from place to place, David Pellow writes that there are “common realities they all share in their experience of colonization that make it possible to generalize an Indigenous methodology while recognizing specific, localized conditions”. Even abstract ideas like the right to a clean environment, a human right according to the United Nations, contradicts Indigenous peoples understanding of environmental justice as it reflects the commodification of land when seen in light of property values.
Environmentalism of the poor
Joan Martinez-Alier's influential concept of the environmentalism of the poor highlights the ways in which marginalized communities, particularly those in the Global South, are disproportionately affected by environmental degradation and the importance of including their perspectives and needs in environmental decision-making. Martinez-Alier's work also introduces the concept of "ecological distribution conflicts," which are conflicts over access to and control of natural resources and the environmental impacts that result from their use, and which are often rooted in social and economic inequalities.
Slow violence
The term “slow violence” was coined by author Rob Nixon in his 2011 book Slow Violence and the Environmentalism of the Poor. Slow violence is defined as “violence that occurs gradually and out of sight, a violence of delayed destruction that is dispersed across time and space, an attritional violence that is typically not viewed as violence at all”. Examples include the effects of climate change, toxic drift, deforestation, oil spills, and the environmental aftermath of war.
Slow violence exacerbates the vulnerability of ecosystems and of people who are poor, disempowered, and often involuntarily displaced, while fueling social conflicts that arise from desperation. Environmental justice as a social movement addresses environmental issues that may be defined as slow violence and otherwise may not be addressed by legislative bodies.
Critical environmental justice
Drawing on concepts of anarchism, posthumanism, critical theory, and intersectional feminism, author David Naguib Pellow created the concept of Critical Environmental Justice (CEJ) in his work What is Critical Environmental Justice. Critical EJ is a perspective intended to address a number of limitations and tensions within EJ Studies. Critical EJ calls for scholarship that builds on research in environmental justice studies by questioning assumptions and gaps in earlier work in the field, embracing greater interdisciplinary, and moving towards methodologies and epistemologies including and beyond the social sciences. Critical EJ scholars believe that since multiple forms of inequality drive and characterize the experience of environmental injustice, the EJ field would benefit from expanding in that direction.
Differentiation between conventional environmental studies and Critical EJ studies is done through four distinctive "pillars". These include, in David Pellow's writing: (1) questions concerning the degree to which scholars should place emphasis on one or more social categories of difference (e.g., race, class, gender, sexuality, species, etc.) versus a focus on multiple forms of inequality; (2) the extent to which scholars studying EJ issues should focus on single-scale versus multi-scalar analyses of the causes, consequences, and possible resolutions of EJ struggles; (3) the degree to which various forms of social inequality and power—including state power—are viewed as entrenched and embedded in society; and (4) the largely unexamined question of the expendability of human and non-human populations facing socioecological threats from states, industries, and other political economic forces.
In his 2017 publication What is Critical Environmental Justice, David Pellow writes as an example of the four pillars working in-tandem:Where we find rivers dammed for hydropower plants we also tend to find indigenous peoples and fisherfolk, as well as other working people, whose livelihoods and health are harmed as a result; when sea life suffers from exposure to toxins such as mercury, we find that human beings also endure the effects of mercury when they consume those animals; and the intersecting character of multiple forms of inequality is revealed when nuclear radiation or climate change affects all species and humans across all social class levels, racial/ethnic groups, genders, abilities, and ages. David Pellow applies his concept of Critical EJ towards modern-day movements in his publication Toward A Critical Environmental Justice Studies, in which he applied the aforementioned pillars towards the Black Lives Matter movement and the problem of state violence. Pellow argues that within conventional studies, “the Black Lives Matter movement and the struggle against environmental racism … is a connection that many scholars might not make at first glance because police brutality and environmental politics would appear to be only tangentially related.” Following his four pillars of Critical EJ, his ties the concept to the Black Lives Matter movement and associated movements, demonstrating:(1) how attention to multiple categories of difference and inequality (including more-than-human species and the built environment); (2) an emphasis on the role of scale as a way of understanding the violence of racism and the promise of resistance movements; (3) a focus on linking the entrenched character of social inequalities with transformative, anti-authoritarian and anarchist perspectives; (4) and an application of the concepts of racial and socioecological indispensability can produce an enriched account of that movement's core concerns, its limitations, and its possibilities.
First Pillar: Intersectionality
The first pillar of Critical EJ Studies involves the recognition that social inequality and oppression in all forms intersect, and that actors in the more-than-human world are subjects of oppression and frequently agents of social change. Developed by Kimberlé Crenshaw in 1989, intersectionality theory states that individuals exist in a crossroads of all their identities, with privilege and marginalization in the intersection between their class, race, gender, sexuality, queerness, cis- or transness, ethnicity, ability, and other facts of identity. As David Nibert and Michael Fox put it in the context of injustice, “The oppression of various devalued groups in human societies is not independent and unrelated; rather, the arrangements that lead to various forms of oppression are integrated in such a way that the exploitation of one group frequently augments and compounds the mistreatment of others.” Thus, Critical EJ views racism, heteropatriarchy, classism ,nativism, ableism, ageism, speciesism (the belief that one species is superior to another), and other forms of inequality as intersecting axes of domination and control.
The organization Intersectional Environmentalism, founded by Leah Thomas in 2020, builds from this theory to argue that intersectional environmentalism means that “social [and] environmental justice are intertwined and environmental advocacy that disregards this connection is harmful and incomplete.”
Second Pillar: Scale
The second pillar of Critical EJ is a focus on the role of scale in the production and possible resolution of environmental injustices. Critical EJ embraces multi-scalar methodological and theoretical approaches order to better comprehend the complex spatial and temporal causes, consequences, and possible resolutions of EJ struggles. Julie Sze writes, “thinking globally and acting locally also demands that people more fully comprehend the relationship between the local and the global or, in other words, to consider scale”. Scale is deeply racialized, gendered, and classed. While the conclusions of climate scientists are remarkably clear that anthropogenic climate change is occurring at a dramatic pace and with increasing intensity. David Pellow writes in his 2016 publication Toward A Critical Environmental Justice Studies that “this is also happening unevenly, with people of color, the poor, indigenous peoples, peoples of the global South, and women suffering the most.”
Pellow further contextualizes scale through temporal dimensions. For instance, how does the emergence and use of coal-fired power plants and petroleum-based economics develop and change over historical periods, and in turn unveiling the social causes of our ecological crises. Pellow observes in his 2017 publication What is Critical Environmental Justice that while “a molecule of carbon dioxide or nitrous oxide can occur in an instant, … it remains in the atmosphere for more than a century, so the decisions we make at one point in time can have dramatic ramifications for generations to come”. Pollution does not stay where it starts, and so consideration must be taken as to the scale of an issue rather than solely its effects.
Third Pillar: Social Inequalities
The third pillar of Critical EJ is the view that social inequalities - from racism to speciesism - are deeply embedded in society and reinforced by state power, and therefore the current social order stands as a fundamental obstacle to social and environmental justice. Pellow argues in his 2017 publication What is Critical Environmental Justice that social change movements may be better off thinking and acting beyond the state and capital as targets of reform and/or as reliable partners. Furthermore, that scholars and activists are not asking how they might build environmentally resilient communities that exist beyond the state, but rather how they might do so with a different model of state intervention.SOURCE Pellow believes that by building and supporting strongly democratic practices, relationships, and institutions, movements for social change will become less dependent upon the state, while any elements of the state they do work through may become more robustly democratic.
He contextualizes this pillar with activist the anarchist-inspired Common Ground Collective, which was co-created by Scott Crow to provide services for survivors of Hurricane Katrina on the Gulf Coast in 2005. Crow gave insight as to what change outside of state power looks like, telling Pellow: We did service work, but it was a revolutionary analysis and practice. We created a horizontal organization that defied the state and did our work in spite of the state … not only did we feed people and give them aid and hygiene kits and things like that, but we also stopped housing from being bulldozed, we cut the locks on schools when they said schools couldn't be opened, and we cleaned the schools out because the students and the teachers wanted that to happen. And we didn't do a one size fits all like the Red Cross would do – we asked the communities, every community we went into, we asked multiple people, the street sex workers, the gangsters, the church leaders, everybody, we talked to them: what can we do to help your neighborhood, to help your community, to help you? And that made us different because for me, it's the overlay of anarchism. Instead of having one franchise thing, you just have concepts, and you just pick the components that match the needs of the people there.
Fourth Pillar: Indispensability
The fourth pillar of Critical EJ centers on a concept David Pellow calls “Indispensability”. Joen Márquez introduces the concept of “racial expendability” in his book Black and Brown Solidarity, in which he argues that “black and brown bodies are, in the eyes of the state and its constituent legal system, generally viewed as criminal, deficient, threatening, and deserving of violent discipline and even obliteration.” Critical EJ builds on this work by countering the ideology of white supremacy and human dominionism, and articulating the perspective that excluded, marginalized, and other populations, beings, and things - both human and nonhuman - must be viewed not as expensable but rather an indispensable to our collective futures. Pellow uses racial indispensability when referring to people of color and socioecological indispensability when referring to broader communities within and across the human/nonhuman divide and their relationships to one another. Pellow expands writing in Toward A Critical Environmental Justice Studies that “racial indispensability is intended to challenge the logic of racial expendability and is the idea that institutions, policies, and practices that support and perpetrate anti-Black racism suffer from the flawed assumption that the future of African Americans is somehow de-linked from the future of White communities.”
History
Traces of environmental injustices span millennia of unrecorded history. Indigenous peoples experienced environmental devastation of a genocidal kind before federal recognition. Origins of the environmental justice movement can be traced to the Indigenous Environmental Movement, which has involved Indigenous populations fighting against displacement and assimilation for sovereignty and land rights for hundreds of years.
The terms 'environmental justice’ and ‘environmental’ racism’ did not enter the common vernacular until residents of Warren County, North Carolina protested against a landfill designed to accept polychlorinated biphenyls in the 1982 PCB protests.
Thirty-thousand gallons of PCB fluid lined 270 miles of roadway in fourteen North Carolina Counties, and the state announced that a landfill would be built rather than undergoing permanent detoxification. Warren County was chosen, the poorest county in the state with a per capita income of around $5,000 in 1980[1], and the site was set for the predominantly Black community of Afton. Its residents protested for six-weeks, leading to over 500 arrests.
That the protests in Warren County were led by civilians led to the basis of future and modern-day environmental, grassroots organizations fighting for environmental justice. Deborah Ferruccio, a contributor to the protest, explained in an interview with The Warren Record that those present were ordinary people. Her husband Ken Ferruccio learned of the PCB dumping after reading newspapers meant for their garden mulch, and days later he and Rev. Leon White led a “humane blockade” to prevent trucks from arriving at the landfill. After being arrested for the demonstration, Furriccio continued his defiance against the county by refusing to post bail and going on a nineteen-day hunger strike.
Rev. Benjamin Chavis was serving for the United Church of Christ (UCC) Commission for Racial Justice when he was sent to Warren County for the protests. Chavis was among the 500 arrested for taking part in the nonviolent protests and is credited with having coined the term “environmental racism” while in the Warren County jail. His involvement, alongside Rev. Leon White, who also served for the UCC, laid the foundation for more activism and consciousness-raising. Chavis would later recall in a New Yorker's article titled “Fighting Environmental Racism in North Carolina” that while “Warren County made headlines … [he] knew in the eighties you couldn't just say there was discrimination. You had to prove it.” Fighting for change, not recognition, is an additional factor of environmental justice as a social movement.
In response to the Warren County Protests, two cross-sectional studies were conducted to determine the demographics of those exposed to uncontrolled toxic waste sites and commercial hazardous waste facilities. The United Church of Christ's Commission for Racial Justice studied the placement of hazardous waste facilities in the US and found that race was the most important factor predicting placement of these facilities. These studies were followed by widespread objections and lawsuits against hazardous waste disposal in poor, generally Black, communities. The mainstream environmental movement began to be criticized for its predominately white affluent leadership, emphasis on conservation, and failure to address social equity concerns.
The EPA established the Environmental Equity Work Group (EEWG) in 1990 in response to additional findings by social scientists that “racial minority and low-income populations bear a higher environmental risk burden than the general population’ and that the EPA's inspections failed to adequately protect low-income communities of color”. In 1992, the EPA published Environmental Equity: Reducing Risks for All Communities - the first time the agency embarked on a systematic examination of environmental risks to communities of color. This acted as their direction of addressing environmental justice.
In 1993 the EPA founded the National Environmental Justice Advisory Council (NEJAC). In 1994 the office's name was changed to the Office of Environmental Justice as a result of public criticism on the difference between equity and justice. SOURCE That same year, President Bill Clinton issued Executive Order 12898, which created the Interagency Working Group on Environmental Justice. The working group sought to address environmental justice in minority populations and low-income populations. David Pellow writes that the executive order “remains the cornerstone of environmental justice regulation in the US, with the EPA as its ventral arbiter”.
Emergence of global movement
Throughout the 1970s and 1980s, grassroots movements and environmental organizations advocated for regulations that increased the costs of hazardous waste disposal in the US and other industrialized nations. However, this led to a surge in exports of hazardous waste to the Global South during the 1980s and 1990s. This global environmental injustice, including the disposal of toxic waste, land appropriation, and resource extraction, sparked the formation of the global environmental justice movement.
Environmental justice as an international subject commenced at the First National People of Color Environmental Leadership Summit in 1991, held in Washington, DC. The four-day summit was sponsored by the United Church of Christ's Commission for Racial Justice. With around 1,100 persons in attendance, representation included all 50 states as well as Puerto Rico, Brazil, Chile, Mexico, Ghana, Liberia, Nigeria, and the Marshall Islands. The summit broadened the environmental justice movement beyond its anti-toxins focus to include issues of public health, worker safety, land use, transportation, housing, resource allocation, and community empowerment. The summit adopted 17 Principles of Environmental Justice, which were later disseminated at the 1992 Earth Summit in Rio, Brazil. The 17 Principles have a likeness in the Rio Declaration on Environment and Development.
In the summer of 2002, a coalition of non-governmental organizations met in Bali to prepare final negotiations for the 2002 Earth Summit. Organizations included CorpWatch, World Rainforest Movement, Friends of the Earth International, the Third World Network, and the Indigenous Environmental Network. They sought to articulate the concept of climate justice. During their time together, the organizations codified the Bali Principles of Climate Justice, a 27-point program identifying and organizing the climate justice movement. Meena Raman, Head of Programs at the Third World Network, explained that in their writing they “drew heavily on the concept of environmental justice, with a significant contribution from movements in the United States, and recognized that economic inequality, ethnicity, and geography played roles in determining who bore the brunt of environmental pollution”. At the 2007 United Nations Climate Conference, or COP13, in Bali, representatives from the Global South and low-income communities from the North created a coalition titled “Climate Justice Now!”. CJN! Issued a series of “genuine solutions” that echoed the Bali Principles.
Initially, the environmental justice movement focused on addressing toxic hazards and injustices faced by marginalized racial groups within affluent nations. However, during the 1991 Leadership Summit, its scope broadened to encompass public health, worker safety, land use, transportation, and other issues. Over time, the movement expanded further to include considerations of gender, international injustices, and intra-group disparities among disadvantaged populations.
Scope
Environmental justice has evolved into a comprehensive global movement, introducing numerous concepts to political ecology, including ecological debt, environmental racism, climate justice, food sovereignty, corporate accountability, ecocide, sacrifice zones, and environmentalism of the poor. It aims to augment human rights law, which traditionally overlooked the relationship between the environment and human rights. Despite attempts to integrate environmental protection into human rights law, challenges persist, particularly concerning climate justice.
Scholars such as Kyle Powys Whyte and Dina Gilio-Whitaker have extended the discourse on environmental justice concerning Indigenous peoples and settler-colonialism. Gilio-Whitaker critiques distributive justice, which assumes a capitalistic commodification of land inconsistent with Indigenous worldviews. Whyte explores environmental justice within the context of colonialism's catastrophic environmental impacts on Indigenous peoples' traditional livelihoods and identities.
Environmental discrimination and conflict
The environmental justice movement seeks to address environmental discrimination and environmental racism associated with hazardous waste disposal, resource extraction, land appropriation, and other activities. This environmental discrimination results in the loss of land-based traditions and economies, armed violence (especially against women and indigenous people) environmental degradation, and environmental conflict. The global environmental justice movement arises from these local place-based conflicts in which local environmental defenders frequently confront multi-national corporations. Local outcomes of these conflicts are increasingly influenced by trans-national environmental justice networks.
There are many divisions along which an unjust distribution of environmental burdens may fall. Within the US, race is the most important determinant of environmental injustice. In other countries, poverty or caste (India) are important indicators. Tribal affiliation is also important in some countries. Environmental justice scholars Laura Pulido and David Pellow argue that recognizing environmental racism, as an element stemming from the entrenched legacies of racial capitalism, is crucial to the movement, with white supremacy continuing to shape human relationships with nature and labor.
Environmental racism
Environmental racism is a pervasive and complex issue that affects communities all over the world. It is a form of systemic discrimination that is grounded in the intersection of race, class, and environmental factors.
At its core, environmental racism refers to the disproportionate exposure of certain communities, mostly those that are marginalized, to environmental hazards such as pollution, toxic waste, and other environmental risks. The Low country Alliance for Model Communities (LAMC) combats environmental racism by empowering marginalized neighborhoods in North Charleston, South Carolina, using community-based research and collaborative problem-solving to identify solutions to health and environmental disparities. These communities are often located near industrial sites, waste facilities, and other sources of pollution that can have serious health impacts.
Environmental racism has a long and troubling history, with many examples dating back to the early 20th century. For instance, the practice of "redlining" in the US, which involved denying loans and insurance to communities of colour, often led to these communities being located in areas with high levels of pollution and environmental hazards.
Today, environmental racism continues to be a significant environmental justice issue, with many low-income communities and communities of colour facing disproportionate exposure to pollution and other environmental risks. This can have serious consequences for the health and well-being of these communities, leading to higher rates of asthma, cancer, and other illnesses.
Addressing environmental racism requires a multifaceted approach that tackles the underlying social, economic, and political factors that contribute to its persistence. More particularly, environmental justice scholars from Latin America and elsewhere advocate to understand this issue through the lens of decolonisation. The latter underlies the fact that environmental racism emanates from the colonial projects of the West and its current reproduction of colonial dynamics.
Hazardous waste
As environmental justice groups have grown more successful in developed countries such as the United States, the burdens of global production have been shifted to the Global South where less-strict regulations make waste disposal cheaper. Export of toxic waste from the US escalated throughout the 1980s and 1990s. Many impacted countries do not have adequate disposal systems for this waste, and impacted communities are not informed about the hazards they are being exposed to.
The Khian Sea waste disposal incident was a notable example of environmental justice issues arising from international movement of toxic waste. Contractors disposing of ash from waste incinerators in Philadelphia, Pennsylvania illegally dumped the waste on a beach in Haiti after several other countries refused to accept it. After more than ten years of debate, the waste was eventually returned to Pennsylvania. The incident contributed to the creation of the Basel Convention that regulates international movement of toxic waste.
Land appropriation
Countries in the Global South disproportionately bear the environmental burden of global production and the costs of over-consumption in Western societies. This burden is exacerbated by changes in land use that shift vast tracts of land away from family and subsistence farming toward multi-national investments in land speculation, agriculture, mining, or conservation. Land grabs in the Global South are engendered by neoliberal ideology and differences in legal frameworks, land prices, and regulatory practices that make countries in the Global South attractive to foreign investments. These land grabs endanger indigenous livelihoods and continuity of social, cultural, and spiritual practices. Resistance to land appropriation through transformative social action is also made difficult by pre-existing social inequity and deprivation; impacted communities are often already struggling just to meet their basic needs.
Water
Access to clean water is an indispensable aspect of human life, yet it remains very unequal, disproportionately affecting marginalized communities globally. The burden of water scarcity is particularly noticeable in impoverished urban settings and remote rural areas where inadequate infrastructure, limited financial resources, and environmental degradation converge to create formidable challenges. Marginalized populations, often already grappling with systemic inequalities, encounter heightened vulnerabilities when it comes to securing safe and reliable water sources. Discriminatory practices can further compound these challenges. The ramifications of limited water access are profound, permeating various facets of daily life, including health, education, and overall well-being. Recognizing and addressing these disparities is not only a matter of justice but also crucial for sustainable development. Consequently, there must be efforts towards implementing inclusive water management strategies that prioritize the specific needs of marginalized communities, ensuring equitable access to this fundamental resource and fostering resilience in the face of global water challenges. One way this has been proposed is through Community Based Participatory Development. When this has been applied, as in the case of the Six Nations Indigenous peoples in Canada working with McMaster University researchers, it has shown how community-led sharing and integrating of science and local knowledge can be partnered in response to water quality.
Resource extraction
Resource extraction is a prime example of a tool based on colonial dynamics that engenders environmental racism. Hundreds of studies have shown that marginalized communities, often indigenous communities, are disproportionately burdened by the negative environmental consequences of resource extraction. Communities near valuable natural resources are frequently saddled with a resource curse wherein they bear the environmental costs of extraction and a brief economic boom that leads to economic instability and ultimately poverty. Indigenous communities living near valuable natural resources face even more discrimination, since they are in most cases simply displaced from their home. Power disparities between extraction industries and impacted communities lead to acute procedural injustice in which local communities are unable to meaningfully participate in decisions that will shape their lives.
Studies have also shown that extraction of critical minerals, timber, and petroleum may be associated with armed violence in communities that host mining operations. The government of Canada found that resource extraction leads to missing and murdered indigenous women in communities impacted by mines and infrastructure projects such as pipelines. The Environmental Justice Atlas, that documents conflicts of environmental justice, demonstrates multiple conflicts with high violence on indigenous populations around resource extraction.
Unequal exchange
Unequal exchange is a term used to describe the unequal economic and trade relationship between countries from the Global North and the Global South. The idea is that the exchange of goods and services between these countries is not equal, with Global North countries benefiting more than the others. This occurs for a variety of reasons such as differences in labor costs, technology, and access to resources. Unequal exchange perceives this framework of trade through the lens of decolonisation: colonial power dynamics have led to a trade system where northern countries can trade their knowledge and technology at a very high price against natural resources, materials and labor at a very low price from southern countries. This is kept in place by mechanisms such as enforceable patents, trade regulations and price setting by institutions such as the World Bank or the International Monetary Fund, where northern countries hold most of the voting power. Hence, unequal exchange is a phenomenon that is based on and perpetuates colonial relationships, as it leads to exploitation and enforces existing inequalities between countries of the Global North and Global South. This interdependence also explains the differences in emissions between northern and southern countries: evidently, since northern countries use many resources and materials of the South, they produce and pollute more.
Health impacts of disparate exposure in EJ communities
Environmental justice communities that are disproportionately exposed to chemical pollution, reduced air quality, and contaminated water sources may experience overall reduced health. Poverty in these communities can be a factor that increases their exposure to occupational hazards such as chemicals used in agriculture or industry. When workers leave the work environment they may bring chemicals with them on their clothing, shoes, skin, and hair, creating further impacts on their families, including children. Children in EJ communities are uniquely exposed, because they metabolize and absorb contaminants differently than adults. These children are exposed to a higher level of contaminants throughout their lives, beginning in utero (through the placenta), and are at greater risk for adverse health effects like respiratory conditions, gastrointestinal conditions, and mental conditions.
Fast fashion exposes environmental justice communities to occupational hazards such as poor ventilation that can lead to respiratory problems from inhalation of synthetic particles and cotton dust. Textile dyeing can also expose EJ communities to toxins and heavy metals when untreated wastewater enters water systems used by residents and for livestock. 95% of clothing production takes place in low- or middle-income countries where the workers are under-resourced.
Erasure of women
Though the environmental justice movement seeks to address discrimination, women have historically been discriminated against as the movement evolves from advocacy to institutional change. While grassroots campaigning activities are often dominated by women, gender inequality is more prevalent in institutionalized activities of organizations dominated by salaried professionals. Women have fought back against this trend by establishing their own domestic and international non-governmental organizations, such as the Women's Earth and Climate Action Network (WECAN) and Women's Earth Alliance (WEA).
The US Environmental Protection Agency's definition of environmental injustice does not include gender, instead mentioning environmental injustice to concern race, color, national origin, and income. Gender inequalities in governing bodies have been noted to have an impact on the nature of decisions made, and so consequently federal legislation and discussion surrounding environmental justice often does not include factors of sex. Authors David Pellow and Robert Brulle write in “Environmental justice: human health and environmental inequalities” that environmental injustices “affect human beings unequally along the lines of race, gender, class and nation, so an emphasis on any one of these will dilute the explanatory power of any analytical approach”.
These inequalities have led to the establishment of the Global Gender and Climate Alliance, set up jointly by the United Nations, the IUCN (International Union for Conservation of Nature), and WEDO (Women, Environment and Development Organization). These have all been founded to raise the profile of gender issues in climate change policymaking.
In environmental law
Cost barriers
One of the prominent barriers to minority participation in environmental justice is the initial costs of trying to change the system and prevent companies from dumping their toxic waste and other pollutants in areas with high numbers of minorities living in them. There are massive legal fees involved in fighting for environmental justice and trying to shed environmental racism. For example, in the United Kingdom, there is a rule that the claimant may have to cover the fees of their opponents, which further exacerbates any cost issues, especially with lower-income minority groups; also, the only way for environmental justice groups to hold companies accountable for their pollution and breaking any licensing issues over waste disposal would be to sue the government for not enforcing rules. This would lead to the forbidding legal fees that most could not afford. This can be seen by the fact that out of 210 judicial review cases between 2005 and 2009, 56% did not proceed due to costs.
Relationships to other movements and philosophies
Climate justice
Climate change and climate justice have also been a component when discussing environmental justice and the greater impact it has on environmental justice communities. Air pollution and water pollution are two contributors of climate change that can have detrimental effects such as extreme temperatures, increase in precipitation, and a rise in sea level. Because of this, communities are more vulnerable to events including floods and droughts potentially resulting in food scarcity and an increased exposure to infectious, food-related, and water-related diseases. Currently, without sufficient treatment, more than 80% of all wastewater generated globally is released into the environment. High-income nations treat, on average, 70% of the wastewater they produce, according to UN Water.
It has been projected that climate change will have the greatest impact on vulnerable populations.
Climate justice has been influenced by environmental justice, especially grassroots climate justice.
Ocean justice
The head of "Ocean Collectiv" and "Urban Ocean Lab", marine biologist, Ayana Elizabeth Johnson describes ocean justice as: "where ocean conservation and issues of social equity meet: Who suffers most from flooding and pollution, and who benefits from conservation measures? As sea levels rise and storms intensify, such questions will only grow more urgent, and fairness must be a central consideration as societies figure out how to answer them"
In December 2023 Biden's administration unveiled a whole strategy to improve ocean justice. The main targets of this strategy:
Repair past injustice when people depending on the ocean and contributing very little to environmental destruction, suffered from the impacts of this destruction on the oceans. Those include Indigenous peoples, African Americans, Hispanic and Latino Americans.
Use the knowledge of indigenous people and marine communities in general for restore ocean justice and help ocean conservation.
Environmental groups supported the decision. According to Beth Lowell, the vice president of Oceana (non-profit group): "Offshore drilling, fisheries management and reducing plastic pollution are just a few of the areas where these voices are needed".
In the official document summarizing the new strategy, the administration gave several examples of past implementation of those principles. One of them is Mai Ka Po Mai a strategy for the management of the Papahānaumokuākea Marine National Monument near the Hawaiian Islands conceived after consultations with native communities.
Environmentalism
Relative to general environmentalism, environmental justice is seen as having a greater focus on the lives of everyday people and being more grassroots. Environmental justice advocates have argued that mainstream environmentalist movements have sometimes been racist and elitist.
Reproductive justice
Many participants in the Reproductive Justice Movement see their struggle as linked with those for environmental justice, and vice versa. Loretta Ross describes the reproductive justice framework as addressing "the ability of any woman to determine her own reproductive destiny" and argues this is "linked directly to the conditions in her community – and these conditions are not just a matter of individual choice and access." Such conditions include those central to environmental justice – including the siting of toxic waste and pollution of food, air, and waterways.
Mohawk midwife Katsi Cook founded the Mother's Milk Project in the 1980s to address the toxic contamination of maternal bodies through exposure to fish and water contaminated by a General Motors Superfund site. In underscoring how contamination disproportionately impacted Akwesasne women and their children through gestation and breastfeeding, this project illustrates the intersections between reproductive and environmental justice. Cook explains that, "at the breasts of women flows the relationship of those generations both to society and to the natural world."
Ecofeminism
Ecofeminst find the intersection between environmentalism and feminist philosophy. Ecofeminism is not to be confused with movements or studies on the health impacts of women in the environment. Researcher and author Sarah Buckingham explains that the basis of ecofeminism is rooted in the argument that “women's equality should not be achieved at the expense of worsening the environment, and neither should environmental improvements be gained at the expense of women.” Its origins are drawn in feminist theory, feminist spirituality, animal rights, social ecology, and antinuclear, antimilitarist organizing. On account of its range of intersectionality, ecofeminism has been criticized for its incoherency and lack of potential in addressing climate crisis.
Ecofeminist concerns are taken up by feminist researchers who participate in environmental organizations or contribute to national and international debates. Examples of such include the National Women's Health Network's research around industrial and environmental health; critiques of reproductive technology and genetic engineering by the Feminist Network of Resistance to Reproductive and Genetic Engineering (FINRRAGE); and critiques of environmental approaches to population control by the Committee on Women, Population, and the Environment.
Around the world
Environmental justice campaigns have arisen from local conflicts all over the world. The Environmental Justice Atlas documented 3,100 environmental conflicts worldwide as of April 2020 and emphasised that many more conflicts remained undocumented.
Africa
Democratic Republic of the Congo
Mining for cobalt and copper in the Democratic Republic of the Congo (DRC) has resulted in environmental injustice and numerous environmental conflicts including
Mutanda mine
Kamoto mine
Tilwezembe mine
Conflict minerals mined in the DRC perpetuate armed conflict.
Ethiopia
Mining for gold and other minerals has resulted in environmental injustice and environmental conflict in Ethiopia including
Lega Dembi mine: thousands of people were exposed to mercury by MIDROC corporation, resulting in poisoned food, death of livestock and many miscairrages and birth defects.
Kenticha mine
Kenya
Kenya has, since independence in 1963, focused on environmental protectionism. Environmental activists such as Wangari Maathai stood for and defend natural and environmental resources, often coming into conflict with the Daniel Arap Moi and his government. The country has suffered Environmental issues arising from rapid urbanization especially in Nairobi, where the public space, Uhuru Park, and game parks such as the Nairobi National Park have suffered encroachment to pave way for infrastructural developments like the Standard Gage Railway and the Nairobi Expressway. One of the environmental lawyers, Kariuki Muigua, has championed environmental justice and access to information and legal protection, authoring the Environmental Justice Thesis on Kenya's milestones.
Nigeria
From 1956 to 2006, up to 1.5 million tons of oil were spilled in the Niger Delta, (50 times the volume spilled in the Exxon Valdez disaster). Indigenous people in the region have suffered the loss of their livelihoods as a result of these environmental issues, and they have received no benefits in return for enormous oil revenues extracted from their lands. Environmental conflicts have exacerbated ongoing conflict in the Niger Delta.
Ogoni people, who are indigenous to Nigeria's oil-rich Delta region have protested the disastrous environmental and economic effects of Shell Oil's drilling and denounced human rights abuses by the Nigerian government and by Shell. Their international appeal intensified dramatically after the execution in 1995 of nine Ogoni activists, including Ken Saro-Wiwa, who was a founder of the nonviolent Movement for the Survival of the Ogoni People (MOSOP).
South Africa
Under colonial and apartheid governments in South Africa, thousands of black South Africans were removed from their ancestral lands to make way for game parks. Earthlife Africa was formed in 1988, making it Africa's first environmental justice organisation. In 1992, the Environmental Justice Networking Forum (EJNF), a nationwide umbrella organization designed to coordinate the activities of environmental activists and organizations interested in social and environmental justice, was created. By 1995, the network expanded to include 150 member organizations and by 2000, it included over 600 member organizations.
With the election of the African National Congress (ANC) in 1994, the environmental justice movement gained an ally in government. The ANC noted "poverty and environmental degradation have been closely linked" in South Africa. The ANC made it clear that environmental inequalities and injustices would be addressed as part of the party's post-apartheid reconstruction and development mandate. The new South African Constitution, finalized in 1996, includes a Bill of Rights that grants South Africans the right to an "environment that is not harmful to their health or well-being" and "to have the environment protected, for the benefit of present and future generations through reasonable legislative and other measures that
prevent pollution and ecological degradation;
promote conservation; and
secure ecologically sustainable development and use of natural resources while promoting justifiable economic and social development".
South Africa's mining industry is the largest single producer of solid waste, accounting for about two-thirds of the total waste stream. Tens of thousands of deaths have occurred among mine workers as a result of accidents over the last century. There have been several deaths and debilitating diseases from work-related illnesses like asbestosis. For those who live next to a mine, the quality of air and water is poor. Noise, dust, and dangerous equipment and vehicles can be threats to the safety of those who live next to a mine as well. These communities are often poor and black and have little choice over the placement of a mine near their homes. The National Party introduced a new Minerals Act that began to address environmental considerations by recognizing the health and safety concerns of workers and the need for land rehabilitation during and after mining operations. In 1993, the Act was amended to require each new mine to have an Environmental Management Program Report (EMPR) prepared before breaking ground. These EMPRs were intended to force mining companies to outline all the possible environmental impacts of the particular mining operation and to make provision for environmental management.
In October 1998, the Department of Minerals and Energy released a White Paper entitled A Minerals and Mining Policy for South Africa, which included a section on Environmental Management. The White Paper states "Government, in recognition of the responsibility of the State as custodian of the nation's natural resources, will ensure that the essential development of the country's mineral resources will take place within a framework of sustainable development and in accordance with national environmental policy, norms, and standards". It adds that any environmental policy "must ensure a cost-effective and competitive mining industry."
Asia
Noah Diffenbaugh and Marshall Burke in their study of inequality in Asia demonstrated the interactionalism of economic inequality and global warming. For instance, globalization and industrialization increased the chances of global warming. However, industrialization also allowed wealth inequality to perpetuate. For example, New Delhi is the epicenter of the industrial revolution in the Indian continent, but there is significant wealth disparity. Furthermore, because of global warming, countries like Sweden and Norway can capitalize on warmer temperatures, while most of the world's poorest countries are significantly poorer than they would have been if global warming had not occurred.
China
In China, factories create harmful waste such as nitrogen oxide and sulfur dioxide which cause health risks. Journalist and science writer Fred Pearce notes that in China "most monitoring of urban air still concentrates on one or at most two pollutants, sometimes particulates, sometimes nitrogen oxides or sulfur dioxides or ozone. Similarly, most medical studies of the impacts of these toxins look for links between single pollutants and suspected health effects such as respiratory disease and cardiovascular conditions." The country emits about a third of all the human-made sulfur dioxide (), nitrogen oxides (), and particulates pollution in the world. The Global Burden of Disease Study, an international collaboration, estimates that 1.1 million Chinese die from the effects of this air pollution each year, roughly a third of the global death toll." The economic cost of deaths due to air pollution is estimated at 267 billion yuan (US$38 billion) per year.
Indonesia
Environmental conflicts in Indonesia include:
The Arun gas field where ExxonMobil's development of a natural gas export industry contributed to the insurgency in Aceh in which secessionist fighters led by the Free Aceh Movement attempted to gain independence from the central government which had taken billions in gas revenues from the region without much benefit to the Aceh province. Violence directed toward the gas industry led Exxon to contract with the Indonesian military for protection of the Arun field and subsequent human rights abuses in Aceh.
Malaysia
Environmental justice movements in Malaysia have arisen from conflicts including:
Lynas Advanced Materials Plant: rare earth processing plant producing over a million tonnes of radioactive waste from 2012-2023.
South Korea
Environmental justice movements in South Korea have arisen from conflicts including:
Saemangeum Seawall
Seoul-Incheon canal
Australia
Australia has suffered from a number of environmental injustices, which have usually been caused by polluting corporate projects geared towards extracting natural resources. For example, discriminatory siting of nuclear and hazardous waste facilities. These projects have been detrimental to local climates, biodiversity, and the health of local citizen populations from poorer economic areas. They have also faced little resistance from local and national governments, who tend to cite their ‘economic’ benefits. However, these projects have faced strong resistance from environmental justice organizations, community, and indigenous groups. Australia has a prominent Indigenous population, and they often disproportionately face some of the worst impacts of these projects.
WestConnex Highway Project, Sydney and New South Wales (NSW)
The WestConnex Highway Project emerged as an answer to Sydney's lack of infrastructure to cope as a fast growing city. The highway project is currently under construction, covers 33 km of new and improved highway, and will link up to the city's M4 and M5 highways.The newest WestConnex toll roads opened in 2019. The NSW government believe that the highway is the ‘missing link’ to the city's problem of traffic congestion, and has argued that the project will provide further economic benefits such as job creation.
The WestConnex Action Group (WAG) have said that residents close to the highway have been negatively affected by its high levels of air pollution, caused by an increase in traffic and unventilated smokestacks in its tunnels. Protesters have also argued that the close proximity of the highway will put children especially at risk.
The highway has faced resistance in a variety of forms, including a long-running occupation camp in Sydney Park, as well as confrontations with police and construction workers that have led to arrests. The WAG has set up a damage register for people whose property has been damaged by the highway, in order to document the extent of the damages, and support those who have been affected. The WAG have done this through campaigning for a damage repaid fund, independent damage assessment and potential class action.
Yeelirrie Uranium Mine, Western Australia
The Yeelirrie Uranium Mine was facilitated by Canadian company Cameco. The mine aimed to dig a 9 km open mine pit and destroy 2,400 hectares of traditional lands, including the Seven Sisters Dreaming Songline, important to the Tjiwarl people. The mine has faced strong resistance from the Tjiwral people, especially its women, for over decade.
The mine is the largest uranium deposit in the country, and uses nine million litres of water, whilst generating millions of tonnes of radioactive waste. Around 36 million tonnes of this waste will be produced whilst the mine is operational, which is set to be until 2043.
A group of Tjiwral women took Cameco to court, to initial success. The Environmental Protection Authority (EPA) halted the mine because it was very likely to wipe out several species, including rare stygofauna, the entire western population of a rare saltbush, and harm other wild life like the Malleefowl, Princess parrot and Greater bilby. The state and federal authorities, however, went against the EPA and approved the mine in 2019.
SANTOS Barossa offshore gas in Timor Sea, Northern Territory (NT)
In March 2021, South Australia Northern Territory Oil Search (SANTOS) invested in the Barossa gas field in the Timor Sea, Northern Territory, to great reception from the NT government, saying that it will provide jobs for the local area. The move was condemned by environmental justice organisations, saying that it will have grave impacts on the climate and biodiversity. Crucially, they stressed that the Tiwi people, owners of the local islands, were not adequately consulted, and were worried that any spills would damage local flatback and Olive Ridley turtle populations.
This disregard for the Tiwi people sparked protests from a number of groups, including one in front of the SANTOS Darwin headquarters demanding an end to the Barossa gas project. In September 2021, a coalition of environmental justice organisations from Australia, South Korea and Japan, united under the name Stop Barossa Gas to oppose the project. In March 2022, the Tiwi people filed for a court injunction to stop KEXIM and Korea Trade and Investment Corporation (Korean development finance institutions) funding the project with almost $1bn. The Tiwi people did this on the basis of a lack of consultation from SANTOS, and the detrimental environmental impacts the project will have. In June 2022, the Tiwi people filed another lawsuit for the same reasons, but this time directly against SANTOS.
Europe
The European Environment Agency (EEA) reports that exposure to environmental harms such as pollution is correlated with poverty, and that poorer countries suffer from environmental harms while higher income countries produce the majority of the pollution. Western Europe has more extensive evidence of environmental inequality.
Romani peoples are ethnic minorities that experience environmental discrimination. Discriminatory laws force Romani people In many countries to live in slums or ghettos with poor access to running water and sewage, or where they are exposed to hazardous wastes.
The European Union is trying to strive towards environmental justice by putting into effect declarations that state that all people have a right to a healthy environment. The Stockholm Declaration, the 1987 Brundtland Commission's Report – "Our Common Future", the Rio Declaration, and Article 37 of the Charter of Fundamental Rights of the European Union, all are ways that the Europeans have put acts in place to work toward environmental justice.
Sweden
Sweden became the first country to ban DDT in 1969. In the 1980s, women activists organized around preparing jam made from pesticide-tainted berries, which they offered to the members of parliament. Parliament members refused, and this has often been cited as an example of direct action within ecofeminism.
United Kingdom
Whilst the predominant agenda of the Environmental Justice movement in the United States has been tackling issues of race, inequality, and the environment, environmental justice campaigns around the world have developed and shifted in focus. For example, the EJ movement in the United Kingdom is quite different. It focuses on issues of poverty and the environment, but also tackles issues of health inequalities and social exclusion. A UK-based NGO, named the Environmental Justice Foundation, has sought to make a direct link between the need for environmental security and the defense of basic human rights. They have launched several high-profile campaigns that link environmental problems and social injustices. A campaign against illegal, unreported and unregulated (IUU) fishing highlighted how 'pirate' fisherman are stealing food from local, artisanal fishing communities. They have also launched a campaign exposing the environmental and human rights abuses involved in cotton production in Uzbekistan. Cotton produced in Uzbekistan is often harvested by children for little or no pay. In addition, the mismanagement of water resources for crop irrigation has led to the near eradication of the Aral Sea. The Environmental Justice Foundation has successfully petitioned large retailers such as Wal-mart and Tesco to stop selling Uzbek cotton.
Building of alternatives to climate change
In France, numerous Alternatiba events, or villages of alternatives, are providing hundreds of alternatives to climate change and lack of environmental justice, both in order to raise people's awareness and to stimulate behaviour change. They have been or will be organized in over sixty different French and European cities, such as Bilbao, Brussels, Geneva, Lyon or Paris.
North and Central America
Belize
Environmental justice movements arising from local conflicts in Belize include:
The government of Belize began granting oil concessions without consulting local communities since 2010, with offshore oil drilling being allowed without consultation with local fishermen or the tourism sector, which are the main economic activities in the area, and affecting Mayan and Garifuna communities. Environmental advocacy group, Oceana, collected over 20,000 signatures in 2011 to trigger a national referendum on offshore oil drilling; however, the government of Belize invalidated over 8,000 signatures, preventing the possibility of an official referendum. In response, Oceana and partner organizations organized an unofficial "People's Referendum," which resulted in 90% of Belizeans voting against offshore exploration and drilling. Belize's Supreme Court declared offshore drilling contracts issued by the Government of Belize in 2004 and 2007 invalid in 2013, but the government reconsidered initiating offshore drilling in 2015, with possible new regulations allowing oil and gas exploration in 99% of Belize's territorial waters. In 2022, Oceana began collecting signatures for another moratorium referendum.
Chalillo Dam
Canada
Environmental justice movements arising from local conflicts in Canada include:
Coastal GasLink pipeline
2020 Canadian pipeline and railway protests
Fairy Creek timber blockade
Grassy Narrows road blockade
Grassy Narrows mercury poisoning
Trans Mountain pipeline
Dominican Republic
Environmental justice movements arising from local conflicts in the Dominican Republic:
Pueblo Viejo mine
Guatemala
Environmental justice movements arising from local conflicts in Guatemala include
Escobal mine
El Salvador
Environmental justice movements arising from local conflicts in El Salvador include:
El Dorado Mine, owned by Pacific Rim Mining Corporation
The Canadian company Pacific Rim Mining Corporation operates a gold mine on the site of El Dorado, San Isidro, in the department of Cabañas. The mine has had hugely negative impacts on the local environment, including the reduction of accessibility to fresh water due to the water intensive mining process, as well as the contamination of the local water supply, which negatively affected the health of local citizens and their live stock. Also, Salvadorian investigators found dangerously high levels of arsenic in two rivers close to the mine.
The operations of the mine has caused conflicts, increased divisions in the community, and prompted threats and violence against opposition to the mine. Following the suspension of the project in 2008 due to resistance from local groups, this violence escalated. As of today, at least half a dozen deaths among local group opposing the mine have been related with the presence of Pacific Rim. The strength of opposition to the mine contributed towards a national movement against the project. In 2008 and 2009, both the incumbent and elected Salvadorian presidents agreed publicly to deny the extension of the licence to Pacific Rim to connote its operations. More recently, the new president Sanchéz Cerén stated “mining is not viable in El Salvador.”
Honduras
Honduras has experienced a number of environmental justice struggles, particularly related to the mining, hydroelectric, and logging industries. One of the most high-profile cases was the assassination of Berta Caceres, a Honduran indigenous and environmental rights activist who opposed the construction of the Agua Zarca Dam on the Gualcarque River. Caceres' murder in 2016 sparked widespread outrage and drew international attention to the risks faced by environmental and indigenous activists in Honduras.
Mexico
Environmental justice movements arising from local conflicts in Mexico include
Dolores mine
El Chanate mine
La Revancha mine
Nicaragua
Environmental justice movements arising from local conflicts in Nicaragua include:
Nicaragua Grand Canal
In 2012, the Nicaraguan government approved the construction of the Grand Canal, which will be 286 km long. A large section of the new canal will run through Lake Nicaragua, which is an important source of fresh water for the country. The canal will also have a width of 83 meters, and depth of 27.5 meters, making it suitable for large-range ships. Related infrastructures include two ports, an airport and an oil pipeline.
Opponents to the construction of the canal, such as the Coordinadora de la comunidad negra creole indígena de Bluefields (CCNCB), fear the impacts it will have on the biodiversity, and protected areas like Bosawás and the Bluefields wetlands. Opponents also fear the impacts on the Indigenous and tribal people that the canal would displace, such as the Miskito, Ulwa and Creole. To date, the Nicaraguan government has not made public the results of various viability studies.
Since the approval of the construction of the canal, environmental justice and indigenous groups have presented petitions for review to national courts, as well as one to the International Human Rights Commission. In 2017, these groups suffered a setback, when the National Court rejected the petition to refuse the "Law of the Grand Canal”.
United States
Definitions of environmental inequality typically emphasize either 'disparate exposure' (unequal exposure to environmental harm) or 'discriminatory intent' (often based on race). Disparate exposure has health and social impacts. Poverty and race are associated with environmental injustice. Poor people account for more than 20% of the human health impacts from industrial toxic air releases, compared to 12.9% of the population nationwide. Some studies that test statistically for effects of race and ethnicity, while controlling for income and other factors, suggest racial gaps in exposure that persist across all bands of income.
States may also see placing toxic facilities near poor neighborhoods as preferential from a Cost Benefit Analysis (CBA) perspective. A CBA may favor placing a toxic facility near a city of 20,000 poor people than near a city of 5,000 wealthy people. Terry Bossert of Range Resources reportedly has said that it deliberately locates its operations in poor neighborhoods instead of wealthy areas where residents have more money to challenge its practices. Northern California's East Bay Refinery Corridor is an example of the disparities associated with race and income and proximity to toxic facilities.
In Seattle, Washington, the Duwamish River Community Coalition (DRCC) was formed in 2001 in response to the designation of the Duwamish River as a Superfund site. DRCC works with local communities and both private and public organizations to address the disparate exposure to air and water pollution families of the Duwamish Valley face. Residents of the Duwamish Valley are a population made of primarily South and Central American immigrants of low income, indigenous peoples, and refugees.
African-Americans
African-Americans are affected by a variety of Environmental Justice issues. One notorious example is the "Cancer Alley" region of Louisiana. This 85-mile stretch of the Mississippi River between Baton Rouge and New Orleans is home to 125 companies that produce one quarter of the petrochemical products manufactured in the United States. The nickname was given due to the high rates of residents diagnosed with cancer compared to the United States average. The United States Commission on Civil Rights has concluded that the African-American community has been disproportionately affected by Cancer Alley as a result of Louisiana's current state and local permit system for hazardous facilities, as well as their low socio-economic status and limited political influence. Another incidence of long-term environmental injustice occurred in the "West Grove" community of Miami, Florida. From 1925 to 1970, the predominately poor, African American residents of the "West Grove" endured the negative effects of exposure to carcinogenic emissions and toxic waste discharge from a large trash incinerator called Old Smokey. Despite official acknowledgement as a public nuisance, the incinerator project was expanded in 1961. It was not until the surrounding, predominantly white neighborhoods began to experience the negative impacts from Old Smokey that the legal battle began to close the incinerator.
More so, many African-American residents have experienced missed or overlooked health issues that were cause by the environmental disparity of their communities. Unfortunately, many of these complications were overlooked by the healthcare industry and comprised the health of those struggling with respiratory and heart problems. The American Heart Association has compiled data analysis that shows the relationship between air pollution exposure and cardiovascular illness and death.
Indigenous Groups
Indigenous groups are often the victims of environmental injustices. Native Americans have suffered abuses related to uranium mining in the American West. Churchrock, New Mexico, in Navajo territory was home to the longest continuous uranium mining in any Navajo land. From 1954 until 1968, the tribe leased land to mining companies who did not obtain consent from Navajo families or report any consequences of their activities. Not only did the miners significantly deplete the limited water supply, but they also contaminated what was left of the Navajo water supply with uranium. Kerr-McGee and United Nuclear Corporation, the two largest mining companies, argued that the Federal Water Pollution Control Act did not apply to them, and maintained that Native American land is not subject to environmental protections. The courts did not force them to comply with US clean water regulations until 1980.
The Inuit community in northern Quebec have faced disproportionate exposure to persistent organic pollutants (POPs) including dioxins and polychlorinated biphenyls (PCBs). Some of these pollutants may include pesticides used decades before in the United States. PCBs bioaccumulate and biomagnify within the fatty tissues of organisms, so the traditional high-fat sea animal diet of the Inuit has posed significant health impacts to both adults and unborn infants. Although the production of PCBs was banned internationally in 2001 by the Stockholm Convention on Persistent Organic Pollutants, they can exist in the environment and biosphere for decades or longer. They pose a significant risk to newborns due to intrauterine exposure and concentration within breast milk.
Latinos
The most common example of environmental injustice among Latinos is the exposure to pesticides faced by farmworkers. After DDT and other chlorinated hydrocarbon pesticides were banned in the United States in 1972, farmers began using more acutely toxic organophosphate pesticides such as parathion. A large portion of farmworkers in the US are working as undocumented immigrants, and as a result of their political disadvantage, are not able to protest against regular exposure to pesticides or benefit from the protections of Federal laws. Exposure to chemical pesticides in the cotton industry also affects farmers in India and Uzbekistan. Banned throughout much of the rest of the world because of the potential threat to human health and the natural environment, Endosulfan is a highly toxic chemical, the safe use of which cannot be guaranteed in the many developing countries it is used in. Endosulfan, like DDT, is an organochlorine and persists in the environment long after it has killed the target pests, leaving a deadly legacy for people and wildlife.
Residents of cities along the US-Mexico border are also affected. Maquiladoras are assembly plants operated by American, Japanese, and other foreign countries, located along the US-Mexico border. The maquiladoras use cheap Mexican labor to assemble imported components and raw material, and then transport finished products back to the United States. Much of the waste ends up being illegally dumped in sewers, ditches, or in the desert. Along the Lower Rio Grande Valley, maquiladoras dump their toxic wastes into the river from which 95 percent of residents obtain their drinking water. In the border cities of Brownsville, Texas, and Matamoros, Mexico, the rate of anencephaly (babies born without brains) is four times the national average.
Youth
Held v. Montana was the first state constitutional law climate lawsuit to go to trial in the United States, on June 12, 2023. The case was filed in March 2020 by sixteen youth residents of Montana, then aged 2 through 18, who argued that the state's support of the fossil fuel industry had worsened the effects of climate change on their lives, thus denying their right to a "clean and healthful environment in Montana for present and future generations":Art. IX, § 1 as required by the Constitution of Montana. On August 14, 2023, the trial court judge ruled in the youth plaintiffs' favor, though the state indicated it would appeal the decision. Montana's Supreme Court heard oral arguments on July 10, 2024, its seven justices taking the case under advisement. On December 18, 2024, the Montana Supreme Court upheld the county court ruling.
South America
Environmental justice struggles have been a significant feature of social and political movements in South America, where communities have faced the impacts of environmental degradation and resource extraction for decades. In particular, mining in South America has led to conflicts between mining companies, governments, and local communities over issues such as land rights, water use, and pollution. Indigenous peoples in particular have been disproportionately affected by mining, with many communities experiencing displacement, loss of traditional livelihoods, and negative health impacts from exposure to toxic chemicals and pollution. A report by Global Witness identifies South America as the most dangerous region in the world for environmental activists, with at least 98 people killed in 2019.
Argentina
Environmental justice movements arising from local conflicts in Argentina include
Bajo de la Alumbrera mine, Catamarca, Argentina: The Bajo de la Alumbrera mine is an open-pit copper and gold mine located in the northwestern province of Catamarca, Argentina. The mining project began in the late 1990s and has since been the center of a significant environmental justice conflict. The mine is operated by Glencore, which owns 50% of the stocks, while Canadian companies Goldcorp and Yamana Gold hold 37.5% and 12.5% respectively. People have raised concerns over the mine's potential environmental impacts, including water pollution, deforestation, and the displacement of indigenous communities. The mine's operators have also faced accusations of human rights violations, including the use of excessive force against protesters and the violation of workers' rights. Despite these concerns, the mine continues to operate, and its expansion plans have been met with significant resistance from local communities and environmental groups. After La Alumbrera started operations, other mining projects were rejected in Catamarca.
Brazil
Environmental justice movements arising from local conflicts in Brazil include
Belo Monte Hydroelectric Dam, Para, Brasil: Belo Monte is a hydroelectric project on the Xingú River in Brazil that began construction in 2011 and was completed in 2019. It is currently the fifth-largest hydroelectric dam in the world, by installed capacity. It is owned by a consortium called Norte Energia, mostly owned by the government and funded primarily by BNDES, with mining giant Vale owning around 5% of it. The project is the largest infrastructure complex of the Brazilian government's plan to build over 60 large dams in the Amazon Basin over the next 20 years, which has received numerous criticisms and open resistance from organizations, public opinion, and inhabitants of the region. Its construction has been highly conflictive, having been opposed by indigenous peoples, who were not consulted before the authorization of construction. The project has been criticized for lacking environmental impact assessments prior to the start of the works. The Belo Monte Dam has diverted the flow of the Xingu, devastating an extensive area of the rainforest, affecting over 50,000 people and displacing over 20,000. The dam threatens the survival of indigenous tribes that depend on the river.
Ecuador
Notable environmental justice movements in Ecuador have arisen from several local conflicts:
Chevron Texaco's oil operations in the Lago Agro oil field resulted in spillage of seventeen million gallons of crude oil into local water supplies between 1967 and 1989. They also dumped over 19 billion gallons of toxic wastewater into unlined open pits and regional rivers. Represented by US lawyer Steven Donziger, Indigenous people fought Chevron in US and Ecuadorian courts for decades in attempts to recover damages.
The Yasuni-ITT Initiative attempted to prevent oil extraction from Yasuni National Park in 2007, but failed and drilling began in 2016.
Peru
Notable environmental justice conflicts in Peru include
Las Bambas copper mine
Yanacocha gold mine
Green space disparities in Lima which has led to higher environmental risks in coastal desert communities compared to wealthier ones
In late March, 2024, the Inter-American Court of Human Rights, based in Costa Rica, ruled that the government of Peru is liable for physical and mental harm to people caused by a metallurgical facility's pollution, and ordered the government to provide free medical care and monetary compensation to victims.
Transnational Movement Networks
Many of the Environmental Justice Networks that began in the United States expanded their horizons to include many other countries and became Transnational Networks for Environmental Justice. These networks work to bring Environmental Justice to all parts of the world and protect all citizens of the world to reduce the environmental injustice happening all over the world. Listed below are some of the major Transnational Social Movement Organizations.
Amazon Watch - organization that campaigns for the protection of the rainforest, and the rights of Indigenous peoples in the Amazon Basin in Ecuador, Peru, Colombia, and Brazil.
Basel Action Network – works to end toxic waste dumping in poor undeveloped countries from the rich developed countries.
—a network of activist-researchers that document environmental justice issues around the world.
Environmental Justice Organisations, Liabilities and Trade (EJOLT) is a multinational project supported by the European Commission. Civil society organizations and universities from 20 countries in Europe, Africa, Latin-America, and Asia are building up case studies, linking organizations worldwide, and making an interactive global map of Environmental Justice.
GAIA (Global Anti-Incinerator Alliance) – works to find different ways to dispose of waste other than incineration. This company has people working in over 77 countries throughout the world.
GR (Global Response) – works to educate activists and the upper working class how to protect human rights and the ecosystem.
Global Witness - an international NGO that investigates and exposes environmental and human rights abuses, corruption, and conflict associated with the exploitation of natural resources.
Greenpeace International – which was the first organization to become the global name of Environmental Justice. Greenpeace works to raise the global consciousness of transnational trade of toxic waste.
Health Care without Harm – works to improve public health by reducing the environmental impacts of the health care industry.
Indigenous Environmental Network - a North American network of indigenous peoples' organizations that work to protect the environment and promote sustainable development.
International Campaign for Responsible Technology – works to promote corporate and government accountability with electronics and how the disposal of technology affect the environment.
International POPs Elimination Network – works to reduce and eventually end the use of persistent organic pollutants (POPs) which are harmful to the environment.
NDN Collective - is an Indigenous-led organization dedicated to building Indigenous, supporting campaigns like ‘Land Back’, which aims to return Indigenous lands back to Indigenous people.
PAN (Pesticide Action Network) – works to replace the use of hazardous pesticides with alternatives that are safe for the environment.
Red Latinoamericana de Mujeres Defensoras de Derechos Ambientales - a regional network that works to promote the rights of women environmental defenders and protect the environment in Latin America.
Global Environmental Activism and Policy
Global environmental inequality is evidence that vulnerable populations are disproportionately victimized by environmental degradation as a result of global capitalism and land exploitation. Yet, studies prove these groups have pioneered the need for intersection between human and environmental rights in activism and policy because of their close proximity to environmental issues. It is important for environmental regulation to acknowledge the value of this global grassroots movement, led by indigenous women and women of the global south, in determining how institutions such as the United Nations can best deliver environmental justice. In recent years, the United Nations' approach to issues concerning environmental health has begun to acknowledge the native practices of indigenous women and advocacy of women in vulnerable positions. Further research by the science community and analysis of environmental issues through a gendered lens are essential next steps for the UN and other governing bodies to curate policy that meets the needs of the women activists leading the environmental justice movement.
Outer space
Over recent years social scientists have begun to view outer space in an environmental conceptual framework. Klinger, an environmental geographer, analyses the environmental features of outer space from the perspective of several schools of geopolitical. From a classical geopolitical approach, for instance, people's exploration of the outer space domain is, in fact, a manifestation of competing and conflicting interests between states, i.e., outer space is an asset used to strengthen and consolidate geopolitical power and has strategic value. From the perspective of environmental geopolitics, the issue of sustainable development has become a consensus politics. Countries thus cede power to international agreements and supranational organizations to manage global environmental issues. Such co-produced practices are followed in the human use of outer space, which means that only powerful nations are capable of reacting to protect the interests of underprivileged countries, so far from there being perfect environmental justice in environmental geopolitics.
Human interaction with outer space is environmentally based since a measurable environmental footprint will be left when modifying the Earth's environment (e.g., local environmental changes from launch sites) to access outer space, developing space-based technologies to study the Earth's environment, exploring space with spacecraft in orbit or by landing on the Moon, etc. Different stakeholders have competing territorial agendas for this vast space; thus, the ownership of these footprints is governed by geopolitical power and relations, which means that human involvement with outer space falls into the field of environmental justice.
Activities on Earth
On Earth, the environmental geopolitics of outer space is directly linked to issues of environmental justice - the launch of spacecraft and the impact of their launch processes on the surrounding environment, and the impact of space-based related technologies and facilities on the development process of human society. As both processes require the support of industry, infrastructure, and networks of information and take place in specific locations, this leads to continuous interaction with local territorial governance.
Launches and infrastructures
Rockets are generally launched in areas where conventional and potentially catastrophic blast damage can be controlled, generally in an open and unoccupied territory. Despite the absence of human life and habitation, other forms of life exist in these open territories, maintaining the local ecological balance and material cycles. Toxic particulate matter from rocket launches can cause localized acid rain, plant and animal mortality, reduced food production, and other hazards.
Moreover, space activities result in environmental injustice on a global scale. Spacecraft are the only contributors to direct human-derived pollution in the stratosphere, which comes mostly from the launch activities of rich economies in the northern hemisphere, while the global north bears more of the environmental consequences.
Environmental injustice is further evidenced by the limited research into the effects on downstream human and non-human communities and the inadequate tracking of pollutants in ecological chains and environments.
Space-based technologies
While space-based technologies have been applied to tracking natural disasters and the spread of pollutants, access to these technologies and the monitoring of data is deeply uneven within and between countries, exacerbating environmental injustice. Further, the use of technology by powerful countries can even lead to the creation of policies and institutions in less privileged nations, changing land-use regimes to favor or disadvantage the survival of certain human groups. For example, in the decades following the publication of the first report on the use of satellite imagery to measure rainforest deforestation in the 1980s, several environmental groups rose to prominence and also influenced changes in domestic policy in Brazil.
| Physical sciences | Earth science basics: General | Earth science |
1863612 | https://en.wikipedia.org/wiki/Normal-form%20game | Normal-form game | In game theory, normal form is a description of a game. Unlike extensive form, normal-form representations are not graphical per se, but rather represent the game by way of a matrix. While this approach can be of greater use in identifying strictly dominated strategies and Nash equilibria, some information is lost as compared to extensive-form representations. The normal-form representation of a game includes all perceptible and conceivable strategies, and their corresponding payoffs, for each player.
In static games of complete, perfect information, a normal-form representation of a game is a specification of players' strategy spaces and payoff functions. A strategy space for a player is the set of all strategies available to that player, whereas a strategy is a complete plan of action for every stage of the game, regardless of whether that stage actually arises in play. A payoff function for a player is a mapping from the cross-product of players' strategy spaces to that player's set of payoffs (normally the set of real numbers, where the number represents a cardinal or ordinal utility—often cardinal in the normal-form representation) of a player, i.e. the payoff function of a player takes as its input a strategy profile (that is a specification of strategies for every player) and yields a representation of payoff as its output.
An example
The matrix provided is a normal-form representation of a game in which players move simultaneously (or at least do not observe the other player's move before making their own) and receive the payoffs as specified for the combinations of actions played. For example, if player 1 plays top and player 2 plays left, player 1 receives 4 and player 2 receives 3. In each cell, the first number represents the payoff to the row player (in this case player 1), and the second number represents the payoff to the column player (in this case player 2).
Other representations
Often, symmetric games (where the payoffs do not depend on which player chooses each action) are represented with only one payoff. This is the payoff for the row player. For example, the payoff matrices on the right and left below represent the same game.
The topological space of games with related payoff matrices can also be mapped, with adjacent games having the most similar matrices. This shows how incremental incentive changes can change the game.
Uses of normal form
Dominated strategies
The payoff matrix facilitates elimination of dominated strategies, and it is usually used to illustrate this concept. For example, in the prisoner's dilemma, we can see that each prisoner can either "cooperate" or "defect". If exactly one prisoner defects, he gets off easily and the other prisoner is locked up for a long time. However, if they both defect, they will both be locked up for a shorter time. One can determine that Cooperate is strictly dominated by Defect. One must compare the first numbers in each column, in this case 0 > −1 and −2 > −5. This shows that no matter what the column player chooses, the row player does better by choosing Defect. Similarly, one compares the second payoff in each row; again 0 > −1 and −2 > −5. This shows that no matter what row does, column does better by choosing Defect. This demonstrates the unique Nash equilibrium of this game is (Defect, Defect).
Sequential games in normal form
These matrices only represent games in which moves are simultaneous (or, more generally, information is imperfect). The above matrix does not represent the game in which player 1 moves first, observed by player 2, and then player 2 moves, because it does not specify each of player 2's strategies in this case. In order to represent this sequential game we must specify all of player 2's actions, even in contingencies that can never arise in the course of the game. In this game, player 2 has actions, as before, Left and Right. Unlike before he has four strategies, contingent on player 1's actions. The strategies are:
Left if player 1 plays Top and Left otherwise
Left if player 1 plays Top and Right otherwise
Right if player 1 plays Top and Left otherwise
Right if player 1 plays Top and Right otherwise
On the right is the normal-form representation of this game.
General formulation
In order for a game to be in normal form, we are provided with the following data:
There is a finite set I of players, each player is denoted by i. Each player i has a finite k number of pure strategies
A is an association of strategies to players, that is an I-tuple
such that
A is a function
whose intended interpretation is the award given to a single player at the outcome of the game. Accordingly, to completely specify a game, the payoff function has to be specified for each player in the player set I= {1, 2, ..., I}.
Definition: A game in normal form is a structure
where:
is a set of players,
is an I-tuple of pure strategy sets, one for each player, and
is an I-tuple of payoff functions.
| Mathematics | Game theory | null |
1864670 | https://en.wikipedia.org/wiki/Marine%20iguana | Marine iguana | The marine iguana (Amblyrhynchus cristatus), also known as the sea iguana, saltwater iguana, or Galápagos marine iguana, is a species of iguana found only on the Galápagos Islands (Ecuador). Unique among modern lizards, it is a marine reptile that has the ability to forage in the sea for algae, which makes up almost all of its diet. Marine iguanas are the only extant lizard that spends time in a marine environment. Large males are able to dive to find this food source, while females and smaller males feed during low tide in the intertidal zone. They mainly live in colonies on rocky shores where they bask after visiting the relatively cold water or intertidal zone, but can also be seen in marshes, mangrove swamps and beaches. Large males defend territories for a short period, but smaller males have other breeding strategies. After mating, the female digs a nest hole in the soil where she lays her eggs, leaving them to hatch on their own a few months later.
Marine iguanas vary in appearance between the different islands and several subspecies are recognized. Although relatively large numbers remain and it is locally abundant, this protected species is considered threatened, primarily from El Niño cycles, introduced predators and chance events like oil spills.
Taxonomy and evolution
Species description and etymology
The marine iguana was first described in 1825 as Amblyrhynchus cristatus by Thomas Bell. He recognized several of its distinctive features, but believed that the specimen he had received was from Mexico, a locality now known to be erroneous.
Its generic name, Amblyrhynchus, is a combination of two Greek words, Ambly- from Amblus (ἀμβλυ) meaning "blunt" and rhynchus (ρυγχος) meaning "snout". Its specific name is the Latin word cristatus meaning "crested," and refers to the low crest of spines along the animal's back.
Amblyrhynchus is a monotypic genus, having only one species, Amblyrhynchus cristatus.
Evolution
Researchers theorize that Galapagos land iguanas (genus Conolophus) and marine iguanas evolved from a common ancestor since arriving on the islands from Central or South America, presumably by rafting. The land and marine iguanas of the Galápagos form a clade, the nearest relatives of which are the Ctenosaura iguanas of Mexico and Central America. Based on a study that relied on mtDNA, the marine iguana was estimated to have diverged from land iguanas some 5-15 million years ago, which is older than any of the extant Galápagos islands. It has therefore traditionally been thought that the ancestral species inhabited parts of the volcanic archipelago that are now submerged. However, a more recent study that included both mtDNA and nDNA indicates that the two split about 4.5 million years ago, which is near the age of the oldest extant Galápagos islands (Española and San Cristóbal).
The marine iguana and Galápagos land iguana remain mutually fertile in spite of being separated by millions of years and assigned to distinct genera. They have been known to hybridize where their ranges overlap, resulting in the so-called hybrid iguana. This is very rare and has only been documented a few times on South Plaza, a tiny island where the usually separated breeding territories and seasons of the two species overlap. The hybrids have intermediate features, stay on land and are infertile.
The different marine iguana populations fall into three main clades: western islands, northeastern islands and southeastern islands. These can be further divided, each subclade generally matching marine iguanas from one or two primary island, except on San Cristóbal where there are two subclades (a northeastern and a southwestern). However, even the oldest divergence between marine iguana populations is quite recent; no more than 230,000 years and likely less than 50,000 years. On occasion one makes it to another island than its home island, resulting in hybridization between different marine iguana populations. There is one confirmed record of a marine iguana outside the Galápagos Islands; in 2014 an individual appeared on Isla de la Plata near the Ecuadorian mainland.
Subspecies
Seven or eight subspecies of the marine iguana, listed alphabetically, have traditionally been recognized:
A. c. albemarlensis Eibl-Eibesfeldt, 1962 – Isabela Island
A. c. ater Gray, 1831 (not always recognized) – Pinzón Island
A. c. cristatus Bell, 1825 – Fernandina Island
A. c. hassi Eibl-Eibesfeldt, 1962 – Santa Cruz Island
A. c. mertensi Eibl-Eibesfeldt, 1962 – San Cristóbal and Santiago Islands
A. c. nanus Garman, 1892 – Genovesa Island
A. c. sielmanni Eibl-Eibesfeldt, 1962 – Pinta Island
A. c. venustissimus Eibl-Eibesfeldt, 1956 – Española Island and adjacent tiny Gardener Island
In 2017, the first comprehensive taxonomic review of the species in more than 50 years came to another result based on morphological and genetic evidence, including recognizing five new subspecies (three of these are small-island populations that were not previously assigned to any subspecies):
A. c. cristatus Bell, 1825 (albermarlensis and ater are junior synonyms) – Isabela Island and Fernandina Island
A. c. godzilla Miralles et al., 2017 – northeastern part of San Cristóbal Island
A. c. hassi Eibl-Eibesfeldt, 1962 – Santa Cruz Island and smaller adjacent islands, such as Baltra Island
A. c. hayampi Miralles et al., 2017 – Marchena Island
A. c. jeffreysi Miralles et al., 2017 – Wolf Island, Darwin Island and Roca Redonda
A. c. mertensi Eibl-Eibesfeldt, 1962 – southwestern part of San Cristóbal Island
A. c. nanus Garman, 1892 – Genovesa Island
A. c. sielmanni Eibl-Eibesfeldt, 1962 – Pinta Island
A. c. trillmichi Miralles et al., 2017 – Santa Fé Island
A. c. venustissimus Eibl-Eibesfeldt, 1956 – Española Island (including adjacent tiny Gardener Island) and Floreana Island
A. c. wikelskii Miralles et al., 2017 – Santiago Island and smaller adjacent islands, such as Rábida Island
Appearance
Early visitors to the Galápagos Islands considered the marine iguanas ugly and disgusting. In 1798, captain James Colnett of the British Royal Navy wrote:
The [i]guanas are small, and of a sooty black, which, if possible, heightens their native ugliness. Indeed, so disgusting is their appearance, that no one on board could be prevailed on, to take them as food.
On his visit to the islands in 1835, despite making extensive observations on the creatures, Charles Darwin was revolted by the animals' appearance, writing:
The black Lava rocks on the beach are frequented by large (2–3 ft []), most disgusting, clumsy Lizards. They are as black as the porous rocks over which they crawl & seek their prey from the Sea. Somebody calls them 'imps of darkness'. They assuredly well-become the land they inhabit.
Marine iguanas have a thickset body and relatively short, robust limbs. Adults have a row of spines extending from the nape, along the back to the tail. Males have longer spines, and larger bony plates on the top of their head than females. There are some differences in the spines depending on the island and they are most elaborate on Fernandina (subspecies cristatus). The scales on the top of the head of the marine iguana are quite conical and pointed. Although less extreme, the same can be seen in the Galápagos land iguanas (genus Conolophus), while these scales at most are slightly domed in other iguanas. Overall the skeleton of the marine iguana shows many similarities with the chuckwallas (genus Sauromalus), but this is an example of homoplasy, as the two are not closely related.
Marine iguanas are not always black; the young have a lighter coloured dorsal stripe, and some adult specimens are grey. Dark tones allow the lizards to rapidly absorb heat to minimize the period of lethargy after emerging from the water. The marine iguana lacks agility on land but is a graceful swimmer. Its laterally flattened tail provides propulsion and the row of spines along the back provide stability, while its long, sharp claws allow it to hold onto rocks in strong currents. Adult males vary in colour with the season, becoming brighter when breeding. There are also major differences in the colour of the adult males depending on the subspecies. Males of the relatively small southern islands of Española, Floreana and Santa Fé (subspecies venustissimus and trillmichi) are the most colourful, with bright pinkish-red and turquoise markings. In comparison, those of the relatively small northern islands of Genovesa, Marchena, Pinta, Wolf and Darwin (jeffreysi, hayampi, sielmanni and nanus) are almost all blackish without contrasting markings. Other subspecies tend to resemble duller versions of venustissimus and trillmichi, or are blackish with markings in pale yellowish, ochre, greenish or grey (sooty to near white). It is suspected that the exact algal diet of each population plays a role in their colour. Females show much less variation between the islands and are typically dark with less contrasting colours than the males. In most places, females do not change colour conspicuously between the breeding and non-breeding season, but at least on Española (subspecies venustissimus) they do assume relatively bright male-like colours when nesting, which is possibly related to their need of defending the nest from other females on an island with few suitable sites.
Size
Marine iguanas typically range from in snout-to-vent length and have a tail length from . There are major geographical differences, as iguanas from large islands tend to grow relatively large as adults, while those from small islands are smaller in size. In one study, the average snout-to-vent length on Wolf and Darwin Islands (subspecies jeffreysi) was about , and those on Genovesa Island (subspecies nanus) were only slightly larger. In comparison, Santa Cruz marine iguanas (hassi) had an average snout-to-vent length of about , and those of Isabela and Fernandina (cristatus) were only slightly smaller. Other subspecies were of intermediate size, in between the small Wolf, Darwin and Genovesa iguanas and the large Santa Cruz, Isabela and Fernandina iguanas. In another study, the largest were from western San Cristóbal Island (mertensi), followed by those from Isabela (cristatus, in part), Floreana (venustissimus, in part), eastern San Cristóbal (godzilla), Fernandina (cristatus, in part) and Santa Cruz (hassi). The smallest by far were from Genovesa (nanus), but this study did not include any marine iguanas from Wolf and Darwin Islands. The remaining island populations were of intermediate size.
The maximum weight of adult males ranges from on southern Isabela to on Genovesa. This difference in body size of marine iguanas between islands is due to the amount of food available, which depends on sea temperature and algae growth.
Marine iguanas are sexually dimorphic with adult males on average being significantly longer and weighing about twice as much as adult females. Additionally, males have broader heads and larger tubercles than females. However, the largest females are only about 20–40% shorter than the largest males. There is a correlation between longevity and body size, particularly for adult males. Large body size in males is selected sexually, but can be detrimental during El Niño events when resources are scarce. This results in large males suffering higher mortality than females and smaller adult males. The mortality rates of marine iguanas are explained through the size difference between the sexes. Some individuals have been shown to shrink in body size by up to 20% during El Niño events and grow back to their previous size when food became available again.
Behavior
Reproduction and life cycle
Female marine iguanas reach sexual maturity at the age of 3–5 years, while males reach sexual maturity at the age of 6–8 years. Sexual maturity is marked by the first steep and abrupt decline in bone growth cycle thickness. Marine iguanas can reach an age of up to 60 years, but average is 12 years or less.
Reproduction in the marine iguana begins during the last part of the cold and dry season, with breeding from December to March and nesting from January to April. The exact timing varies with location, depending on algal growth and the nutrient-rich Cromwell Current. It occurs earliest on islands like Fernandina, Isabela, Santa Fé and Genovesa, and latest on islands like Española. An adult marine iguana, whether male or female, will typically breed every other year, but if there is plenty of food a female may breed each year.
Marine iguanas live in coastal colonies that typically range from 20 to 500 animals, but sometimes more than 1,000. Their bodies often touch each other, but there are no social interactions like grooming, as commonly seen in social mammals and birds. Females are always highly gregarious and males are gregarious outside the breeding season. However, large males defend territories for up to three months during the breeding season, and in this period they sometimes fight other males. A territory can be up to almost , but is usually less than half that size, and can be as small as . A territory is often delimited by geographic features, like boulders or crevices. The territories tend to occur in clusters with several located adjacent to each other, forming a lek, but they can occur singly. Medium-sized males roam areas near the territories of large males or walk along beaches looking for females, while small "sneaky" males may enter the territories of large males. To threaten another male, a male will bob his head, walk on stiff legs, raise the spiny crest along the back and open the mouth to reveal the reddish interior. In most cases the suitor will turn away, but if he responds with the same behavior a fight ensues. During fights they typically do not bite each other, instead thrusting their heads together in an attempt to push the other away. The bony plates on the top of their heads are especially suited for interlocking. Fights between males may last for hours, and are often interrupted by periodic breaks. Once a winner has been established through the headbutting, the loser assumes a submissive position and retreats. In general fights between males are harmless and highly ritualized, but on occasion they will bite and injure each other.
Males are primarily selected by females on the basis of their body size. Females display a stronger preference for mating with bigger males. It is precisely because of body size that reproductive performance increases and "is mediated by higher survival of larger hatchlings from larger females and increased mating success of larger males." Other factors involved in the female's choice of partner are the display frequency by a male (especially head-bobbing) and the quality of a male's territory. Females prefer male territories that include or are near their own resting places, which they choose based on proximity to the sea, access to shade, low midday temperature and the possibility of sunbathing in the afternoon. Males with territories that are near the center of the lek tend to have a greater mating success than males with peripheral territories, but the size of a territory does not affect mating success. Large territorial males that frequently display also emit higher levels of certain acidic compounds (including 11-Eicosenoic acid) from their femoral pores, which may function as pheromones that enhance their chance of attracting females. Females can move freely between different territories, but receive less harassment from opportunistic non-territorial males when inside another male's territory. Medium-sized males attempt to mate by force and small males by stealth and force, but they have a low mating success rate compared to the large males that maintain a territory. During courtship display, a territorial male nods at the female, may open his mouth, and performs a slow sideways walk towards her. If she accepts, the male will mount her while holding her by the neck. A mating lasts no more than 20 minutes, typically 3 to 4 minutes, but it is comparatively rapid in the small "sneaky" males, which easily are overlooked because their size, general morphology and colours are similar to those of a female. This rapid mating is necessary because large males will chase them out of their territory as soon as they are discovered. During each breeding season, a male will mate with many females if given the chance, but the female only mates once. Once a female has mated, she rejects other suitors by nodding her head at them.
Roughly one month after copulation, the female lays between one and six eggs, typically two or three. The leathery white eggs measure about and weigh . This is large for an iguana, and altogether the eggs may weigh up to one-quarter the weight of the female. The nest sites can be as much as inland, but typically are much closer to the coast. They are occasionally as little as inland, although they have to be above the high tide water mark. The nest is deep and dug in sand or volcanic ash by the female. On islands where there are few suitable sites and digging is difficult due to a relatively hard soil and many rocks, the female guards the nest for several days after the eggs have been buried, ensuring that they are not dug up by other nesting females. As in males defending their territory from other males, females defending their nest site from other females begin with a threat display. If this fails to scare the opponent away, the fights between females involve much biting and are less ritualized than the territorial fights between males. Where there are more suitable sites and the soil is loose, females are less likely to fight for a location and do not guard their nest after the eggs have been buried. The eggs hatch after about three to four months. The hatchlings are in snout-to-vent length, and weigh . As soon as they emerge from the nest they run for cover, and begin their trip to locations that provides optimum conditions for feeding, temperature regulation and shelter. Some hatchlings have been recorded moving as far as in two days.
Feeding
The marine iguana forages almost exclusively on red and green algae in the inter- and subtidal zones. At least 10 genera of algae are regularly consumed, including the red algae Centroceras, Gelidium, Grateloupia, Hypnea, Polysiphonia and Pterocladiella. In some populations the green algae Ulva dominates the diet. The algal diet varies depending on algal abundance, individual preferences, foraging behaviour, season and exact island of feeding. Some species with chemical deterrents, such as Bifurcaria, Laurencia and Ochtodes, are actively avoided, but otherwise algal food choice mainly depends on energy content and digestibility. On Santa Cruz Island, 4–5 red algal species are the food of choice. During neap low tides, however, the usually avoided green algae Ulva lobata is eaten more often since the preferred red algae is not easily available. Brown algae have occasionally also been recorded in their diet, but marine iguanas are unable to easily digest these and will starve if it is the only type present. A marine iguana typically eats about dry weight or fresh weight of algae per day. At Punta Espinoza on northeastern Fernandina Island it has been estimated that the almost 1,900 marine iguanas eat about (fresh weight) of algae per year, a rate of consumption that is counterbalanced by the very high growth rate of the algae. They may feed on octopuses, crustaceans, insects (such as grasshoppers and cockroaches), fish carrion, and sea lion feces and afterbirth on rare occasions. The population on North Seymour Island will supplement their diet with land plants, primarily Batis maritima, or other coastal succulents like Sesuvium portulacastrum. These North Seymour iguanas have a higher survival rate during periods where their normal algal food is reduced. However, the hindgut of marine iguanas is specially adapted to algae feeding, likely restricting the possibility of efficiently switching to other plant types. The algae are digested with the help of endosymbiotic bacteria in their gut.
In the first months after hatching, the juveniles mainly feed on feces from larger marine iguanas, gaining the bacteria needed for digesting algae. It has been suggested that young iguanas up to about two years old are unable to swim, but studies have shown that even newly hatched marine iguanas are good swimmers; they just strongly try to avoid entering the water. At about 1–2 years old the young may voluntarily swim in shallow water and tide pools, but they do not dive.
Marine iguanas can dive as deep as , and can spend up to one hour underwater. When diving to or deeper, they regularly remain submerged from 15 to more than 30 minutes. Most dives are much shorter in duration and shallower than . Individuals foraging near-shore, often less than deep, typically only spend about 3 minutes underwater. Only 5% of marine iguanas dive for algae offshore and these individuals are the large males. The minimum size of these divers vary with island and subspecies, ranging from on Genovesa Island (A. c. nanus) to on Fernandina Island (A. c. cristatus). They are slow swimmers, averaging just . The highest recorded speed is only about twice that figure and this can typically only be sustained in bursts that last less than a minute. Most females and smaller males feed on exposed algae in the intertidal zone during low tide, retreating once the water returns and starts washing over them. They often scurry back-and-forth repeatedly, running to a patch of algae to take a few bites and then return fast to higher ground to avoid incoming waves. The separation in feeding behavior is advantageous because the large offshore feeding males experience less competition for food from smaller males and females. A few individuals of intermediate size may use both feeding strategies. In general, each marine iguana has a specific feeding site it returns to day after day. Most feed daily, but large offshore feeding males often only every second or third day. During bad weather with high waves marine iguanas do not feed, sometimes for more than a week. Large males often do not feed for several weeks when maintaining a breeding territory, resulting in them losing up to about one-quarter of their weight. It takes many months for them to return to their original weight. In captivity, individuals have remained strong and active even after fasting for as much as 100 days.
Foraging behavior changes in accordance to the seasons and foraging efficiency increases with temperature. These environmental changes and the ensuing occasional food unavailability have caused marine iguanas to evolve by acquiring efficient methods of foraging in order to maximize their energy intake and body size. During an El Niño cycle in which food diminished for two years, some were found to decrease their length by as much as 20%. When food supply returned to normal, iguana size followed suit. It is speculated that the bones of the iguana actually shorten as shrinkage of connective tissue could only account for a 10% change in length.
Marine iguanas have several adaptions that aid their feeding. Their flattened tail is the primary means of propulsion in the water. When on the surface, they may use their legs for maintaining balance. Although their partially webbed feet often are mentioned, this webbing is very marginal and no greater in extent than in the green iguana, a species that also shares the flattened tail. Marine iguanas have powerful limbs with long, sharp claws for climbing, holding onto rocks and pulling themselves forward when at the sea bottom (adding to the propulsion provided by the tail). They are buoyant and float to the ocean surface if they are not actively swimming or holding on to rocks underwater. However, they have unusually compact (osteosclerose) limb bones compared to the land iguana, especially those from the front limbs, providing ballast to help with diving.
Other adaptions in marine iguanas are blunt heads (flat noses) and sharp teeth allowing them to graze algae off of rocks more easily. Together with a few Ctenosaurus species, it is the only iguana that never has more than three tips (tricuspid) on each tooth. Uniquely, the side-tips on the marine iguana's teeth are quite large, only somewhat smaller than the central tip. It also appears to replace its teeth at a higher rate than other iguanas. As a sea reptile, much salt is ingested. The salt is filtered from their blood and then excreted by specialised cranial exocrine glands at the nostrils, expelled from the body in a process much like sneezing. The marine iguana's cranium has an unusually large nasal cavity compared to other iguanas, which is necessary to accommodate the large salt glands. The head may appear white from encrusted salt.
Mutualism and commensalism with other animals
Marine iguanas have mutualistic and commensal relationships with several other animals. Lava lizards may scurry over marine iguanas when hunting flies, and Darwin's finches, mockingbirds and Sally lightfoot crabs sometimes feed on mites and ticks that they pick off their skin. Marine iguanas typically ignore these visits. When underwater, they are often cleaned by fish, like Pacific sergeant majors that pick off moulting skin. Although there are no apparent benefits to either species, marine iguanas commonly live close together with the much larger Galápagos sea lions. The two species generally ignore each other and an iguana may even crawl over the body of a sea lion.
Thermoregulation
Marine iguanas are unique as they are marine reptiles that forage on inter- and subtidal algae almost exclusively. They forage in the relatively cold waters around the Galápagos Islands, which typically are between at the sea surface. As their preferred body temperature is from and the temperature declines throughout a foraging trip to the sea, sometimes by as much as , the marine iguana has several behavioral adaptations for thermoregulation. At cold temperatures their muscles are less efficient, but their relatively high temperature preference is also related to the optimal temperature for digesting the algal food in their gut.
As an ectothermic animal, the marine iguana can spend only a limited time in cold water diving for algae. Afterwards it basks in the sun to warm up. Until it can do so it is unable to move effectively, making it vulnerable to predation. However, this is counteracted by their highly aggressive nature consisting of biting and expansive bluffs when in this disadvantageous state. Their dark shade aids in heat reabsorption. In colder periods with cloudy weather and much wind, juveniles will stay in the lee of rocks, still gaining the heat from the sun. Adults may move inland to low-lying sites with less wind because of bushes and lava ridges but still exposed to direct sun. When in the water and their temperature is falling, their blood circulation is reduced because of a low heart rate of about 30 beats per minute, allowing them to better conserve their warmth. When on land and heating up, the higher heart rate of about 100 beats per minute aids in spreading the heat throughout the body. To conserve heat during the night, they often sleep closely together in groups that may number up to 50 individuals, while others sleep alone below plants or in crevices.
In general, the time of each foraging trip is directly related to the water temperature; the colder the water the shorter the foraging trip. Additionally, marine iguanas that forage in or near the intertidal zone prefer to do so during low tides, allowing them to remain on land (on rocks exposed by the tide) or return to land faster. Individuals that forage further from the shore by diving are large males, which mainly feed during the hot midday (although it may occur from late morning to early evening), are less affected by the cool water because of their body size and are more efficient swimmers.
Under the tropical sun, overheating can also be a problem. To avoid this, they pant, and adopt a posture where they face the sun and lift their body up, thereby exposing as little as possible of their body to direct sun and allowing cooling air to pass underneath.
Conservation
Status and threats
The marine iguana has a relatively small range and is currently considered vulnerable by the IUCN. Most subpopulations have the same IUCN rating, but those of San Cristóbal, Santiago and Genovesa Islands are considered endangered.
On some shorelines they can be very numerous, with densities as high as 8,000 per kilometer (almost 13,000 per mile), and their biomass compared to the area they occupy may surpass that of any known reptile. However, their distribution is patchy, and colonies are generally found within of the ocean, naturally limiting their range. The total population for the entire archipelago is estimated to be 200,000–300,000 individuals, although this number is labelled with considerable uncertainty. Most subpopulations have not been surveyed in detail because their lifestyle and habitat make it difficult to survey with a high level of accuracy. By far the largest subpopulation—likely including around of all marine iguanas—lives on Fernandina Island, but estimates vary greatly from 15,000 to 120,000 individuals. In contrast, the maximum size of the subpopulations on the islands of San Cristóbal, Darwin and Pinzon is estimated to be 400, 800 and 900, respectively. It is estimated that Marchena Island has 4,000–10,000 marine iguanas, Rabida Island has 1,000–2,000 and Santa Fé Island has 15,000–30,000. Although individuals may on occasion be transferred between islands by ocean currents, marine iguanas are unable to swim between all but the nearest islands in the archipelago because of their slow speed and limited stamina in the relatively cold water.
The periodic El Niño events reduce the cold water needed for algae to grow and this can drastically reduce the marine iguana population, on some islands with as much as 90%. Population recoveries after El Niños are fast; even when reduced by 30–50%, the population is generally able to recover within four years. However, recoveries can be partially impeded by the invasive brown algae Giffordia mitchelliae. When their food algae (red and green algae) disappears during El Niños, the areas may be taken over by this brown algae, causing starvation among the marine iguanas. With global warming, it is expected that El Niño events will be stronger and occur more frequently.
Introduced predators, to which they have little or no protection, include animals such as pigs, dogs, cats and rats. Dogs may take adult marine iguanas, while the others may feed on their young or eggs. This inhibits reproduction and the long-term survival of the species. Introduced predators represent a major problem on the islands of Santa Cruz, Santiago, Isabela, Floreana and San Cristóbal where very few marine iguana hatchlings survive. Marine iguanas are ecologically naïve and have not developed efficient anti-predator strategies against the introduced species. For example, the first dogs were introduced to the Galápagos Islands more than 100 years ago, but marine iguanas have not developed an anti-predator strategy against them. In general, native predators represent a less serious threat to the marine iguana. Natural land predators include the Galápagos hawk, short-eared owl, lava gull, herons and Galápagos racer snakes that may take small marine iguanas. When swimming, marine iguanas are occasionally attacked and eaten by sharks, although the two often behave indifferently to each other, even when close together. Of the native predators, the Galápagos hawk is likely the most important, and it may also take weakened adults (not just young), but this hawk is quite rare with a total population numbering only a few hundred individuals. Marine iguanas have anti-predator strategies that reduce the impact of the Galápagos hawk, including an increased vigilance when hearing the alarm call of the Galápagos mockingbird, another species that is sometimes preyed upon by the hawk.
Marine iguanas can easily be approached by humans as they are very tame and generally make little or no attempt to escape. Individuals that have been caught before are only slightly more wary when again encountering humans. Even when repeatedly caught and deliberately mishandled each time, they do not attempt to bite or lash with their tail in self-defense (although the sharp claws can cause painful scratches when the iguana attempts to gain a grip) and only move a few feet once released, allowing themselves to be caught again with ease. Pathogens introduced to the archipelago by humans pose a serious threat to this species. The marine iguana has evolved over time in an isolated environment and lacks immunity to many pathogens. As a result, the iguanas are at higher risk of contracting infections, contributing to their threatened status.
Occasional oil spills in the region also present a threat. For example, the Santa Fé population was reduced by almost as a result of the MV Jessica oil spill in 2001, and even low-level oiling may kill marine iguanas. It is suspected that the primary cause of death during these events is starvation due to the loss of their sensitive gut bacteria, which they rely on to digest algae.
Protection
The marine iguana is completely protected under the laws of Ecuador, and is listed under CITES Appendix II. Almost all its land range is in the Galápagos National Park (only the 3% human-inhabited sections in the archipelago are excluded) and all its sea range is in the Galápagos Marine Reserve. Certain coastal roads have lower speed limits to reduce the risk posed by cars, especially to the young. There have been attempts of removing introduced predators from some places, but this has not been without problems. Feral dogs mostly feed on large marine iguanas, but also limit the feral cat population. When feral dogs are removed feral cats can thrive and they feed on small marine iguanas. In 2012, the last black rats were eradicated from Pinzón Island.
Studies and research have been done on marine iguanas that can help and promote conservation efforts to preserve the endemic species. Monitoring levels of marine algae, both dimensionally and hormonally, is an effective way to predict the fitness of the marine iguana species. Exposure to tourism affects marine iguanas, and corticosterone levels can predict their survival during El Niño events. Corticosterone levels in species measure the stress that they face in their populations. Marine iguanas show higher stress-induced corticosterone concentrations during famine (El Niño) than feast conditions (La Niña). The levels differ between the islands, and show that survival varies throughout them during an El Niño event. The variable response of corticosterone is one indicator of the general public health of the populations of marine iguanas across the Galápagos Islands, which is a useful factor in the conservation of the species.
Another indicator of fitness is the levels of glucocorticoid. Glucocorticoid release is considered beneficial in helping animals survive stressful conditions, while low glucocorticoid levels are an indicator of poor body condition. Species undergoing a large measure of stress, resulting in elevated glucocorticoid levels can cause complications such as reproduction failure. Human activity has been considered a cause of elevated levels of glucocorticoid in species. Results of a study show that marine iguanas in areas central to tourism are not chronically stressed, but do show lower stress response compared to groups undisturbed by tourism. Tourism, thus, does affect the physiology of marine iguanas. Information of glucocorticoid levels are good monitors in predicting long term consequences of human impact.
Although marine iguanas have been kept in captivity, the specialised diet represents a challenge. They have lived for more than a decade in captivity, but have never bred under such conditions. The development of a captive breeding program (as already exists for the Galápagos land iguana) possibly is a necessity if all the island subspecies are to survive.
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| Biology and health sciences | Iguania | Animals |
1864695 | https://en.wikipedia.org/wiki/Dracaena%20trifasciata | Dracaena trifasciata | Dracaena trifasciata is a species of flowering plant in the family Asparagaceae, native to tropical West Africa from Nigeria east to the Congo. It is most commonly known as the snake plant, Saint George's sword, mother-in-law's tongue, and viper's bowstring hemp, among other names. Until 2017, it was known under the synonym Sansevieria trifasciata. This plant is often kept as a houseplant due to its non-demanding maintenance; they can survive with very little water and sun.
Description
It is an evergreen perennial plant forming dense strands, spreading by way of its creeping rhizome, which is sometimes above ground, sometimes underground. Its stiff leaves grow vertically from a basal rosette. Mature leaves are dark green with light gray-green cross-banding and usually range from long and wide, though it can reach heights above in optimal conditions.
The specific epithet trifasciata means "three bundles".
The plant exchanges oxygen and carbon dioxide using the crassulacean acid metabolism process, which allows them to withstand drought. The microscopic pores on the plant's leaves, called the stomata and used to exchange gases, are opened only at night to prevent water from escaping via evaporation in the hot sun.
To get this plant to go into bloom outside of its natural environment is difficult. Replicating its natural environment is possible. Its flowers vary from greenish white to cream-colored — some are fragrant at night, others not at all — and have a sticky texture.
Dracaena trifasciata is commonly called "mother-in-law's tongue", "Saint George's sword" or "snake plant", because of the shape and sharp margins of its leaves that resemble snakes. It is also known as the "viper's bowstring hemp", because it is one of the sources for plant fibers used to make bowstrings.
Cultivation and uses
Like other members of its genus, D. trifasciata yields bowstring hemp, a strong plant fiber once used to make bowstrings.
It is now used predominantly as an ornamental plant, outdoors in warmer climates, and indoors as a houseplant in cooler climates. It is popular as a houseplant because it is tolerant of low light levels and irregular watering; during winter, it needs only one watering every couple of months. It will rot easily if overwatered. It is commonly recommended to beginners interested in cultivating houseplants for its easy care.
The NASA Clean Air Study found D. trifasciata has the potential to filter indoor air, removing four of the five main toxins involved in the effects of sick building syndrome. However, its rate of filtration is too slow for practical indoor use.
It can be propagated by cuttings or by dividing the rhizome. The first method has the disadvantage that the variegation will be lost.
D. trifasciata is considered by some authorities a potential weed in Australia. It is widely used as an ornamental in the tropics outdoors in pots and garden beds, and in temperate areas as an indoor plant.
The plant contains saponins which are mildly toxic to dogs and cats and can lead to gastrointestinal upset if consumed.
In South Africa, it is used to treat ear infections.
Varieties and cultivars
Numerous cultivars have been developed, many of them for variegated foliage with yellow or silvery-white stripes on the leaf margins. Popular cultivars include 'Compacta', 'Goldiana', 'Hahnii', 'Laurentii', 'Silbersee', and 'Silver Hahnii'. 'Hahnii' was discovered in 1939 by William W. Smith Jr. in the Crescent Nursery Company, New Orleans, Louisiana. The 1941 patent was assigned to Sylvan Frank Hahn of Pittsburgh, Pennsylvania.
The variety D. trifasciata var. laurentii, together with the cultivars 'Bantel's Sensation'
and 'Golden Hahni' have gained the Royal Horticultural Society's Award of Garden Merit.
Non-variegated forms of D. trifasciata are often incorrectly sold as Sansevieria zeylanica, which is a different species that is rarely cultivated.
Cultural significance
In its native range in Africa, Dracaena trifasciata specimens with yellow stripes on the leaf margins are associated with Ọya, the female orisha of storms. In Nigeria, the plant is commonly linked with Ògún, the orisha of war, and is used in rituals to remove the evil eye.
In Brazil, where it is known as espada de São Jorge ("Saint George's sword"), it is grown outside houses to ward off evil that might harm the home (as is Dracaena angolensis, Saint George's spear). The plant plays an important part in the Afro-Brazilian syncretic religion Umbanda, also representing the orisha Ogum (Ògún), as Ogum is syncretized with Saint George. Some yellow-edged varieties of D. trifasciata are called espada de Santa-Bárbara ("Sword of Saint Barbara") and are associated with Iansã, the Umbanda name for Ọya, Saint Barbara's syncretic orisha pair. These types are grown to protect against inclement weather.
This plant is visible on the porch in American Grant Wood's 1930 painting, American Gothic.
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| Biology and health sciences | Asparagales | Plants |
1864889 | https://en.wikipedia.org/wiki/Cosmology | Cosmology | Cosmology () is a branch of physics and metaphysics dealing with the nature of the universe, the cosmos. The term cosmology was first used in English in 1656 in Thomas Blount's Glossographia, and in 1731 taken up in Latin by German philosopher Christian Wolff in Cosmologia Generalis. Religious or mythological cosmology is a body of beliefs based on mythological, religious, and esoteric literature and traditions of creation myths and eschatology. In the science of astronomy, cosmology is concerned with the study of the chronology of the universe.
Physical cosmology is the study of the observable universe's origin, its large-scale structures and dynamics, and the ultimate fate of the universe, including the laws of science that govern these areas. It is investigated by scientists, including astronomers and physicists, as well as philosophers, such as metaphysicians, philosophers of physics, and philosophers of space and time. Because of this shared scope with philosophy, theories in physical cosmology may include both scientific and non-scientific propositions and may depend upon assumptions that cannot be tested. Physical cosmology is a sub-branch of astronomy that is concerned with the universe as a whole. Modern physical cosmology is dominated by the Big Bang Theory which attempts to bring together observational astronomy and particle physics; more specifically, a standard parameterization of the Big Bang with dark matter and dark energy, known as the Lambda-CDM model.
Theoretical astrophysicist David N. Spergel has described cosmology as a "historical science" because "when we look out in space, we look back in time" due to the finite nature of the speed of light.
Disciplines
Physics and astrophysics have played central roles in shaping our understanding of the universe through scientific observation and experiment. Physical cosmology was shaped through both mathematics and observation in an analysis of the whole universe. The universe is generally understood to have begun with the Big Bang, followed almost instantaneously by cosmic inflation, an expansion of space from which the universe is thought to have emerged 13.799 ± 0.021 billion years ago. Cosmogony studies the origin of the universe, and cosmography maps the features of the universe.
In Diderot's Encyclopédie, cosmology is broken down into uranology (the science of the heavens), aerology (the science of the air), geology (the science of the continents), and hydrology (the science of waters).
Metaphysical cosmology has also been described as the placing of humans in the universe in relationship to all other entities. This is exemplified by Marcus Aurelius's observation that a man's place in that relationship: "He who does not know what the world is does not know where he is, and he who does not know for what purpose the world exists, does not know who he is, nor what the world is."
Discoveries
Physical cosmology
Physical cosmology is the branch of physics and astrophysics that deals with the study of the physical origins and evolution of the universe. It also includes the study of the nature of the universe on a large scale. In its earliest form, it was what is now known as "celestial mechanics," the study of the heavens. Greek philosophers Aristarchus of Samos, Aristotle, and Ptolemy proposed different cosmological theories. The geocentric Ptolemaic system was the prevailing theory until the 16th century when Nicolaus Copernicus, and subsequently Johannes Kepler and Galileo Galilei, proposed a heliocentric system. This is one of the most famous examples of epistemological rupture in physical cosmology.
Isaac Newton's Principia Mathematica, published in 1687, was the first description of the law of universal gravitation. It provided a physical mechanism for Kepler's laws and also allowed the anomalies in previous systems, caused by gravitational interaction between the planets, to be resolved. A fundamental difference between Newton's cosmology and those preceding it was the Copernican principle—that the bodies on Earth obey the same physical laws as all celestial bodies. This was a crucial philosophical advance in physical cosmology.
Modern scientific cosmology is widely considered to have begun in 1917 with Albert Einstein's publication of his final modification of general relativity in the paper "Cosmological Considerations of the General Theory of Relativity" (although this paper was not widely available outside of Germany until the end of World War I). General relativity prompted cosmogonists such as Willem de Sitter, Karl Schwarzschild, and Arthur Eddington to explore its astronomical ramifications, which enhanced the ability of astronomers to study very distant objects. Physicists began changing the assumption that the universe was static and unchanging. In 1922, Alexander Friedmann introduced the idea of an expanding universe that contained moving matter.
In parallel to this dynamic approach to cosmology, one long-standing debate about the structure of the cosmos was coming to a climax – the Great Debate (1917 to 1922) – with early cosmologists such as Heber Curtis and Ernst Öpik determining that some nebulae seen in telescopes were separate galaxies far distant from our own. While Heber Curtis argued for the idea that spiral nebulae were star systems in their own right as island universes, Mount Wilson astronomer Harlow Shapley championed the model of a cosmos made up of the Milky Way star system only. This difference of ideas came to a climax with the organization of the Great Debate on 26 April 1920 at the meeting of the U.S. National Academy of Sciences in Washington, D.C. The debate was resolved when Edwin Hubble detected Cepheid Variables in the Andromeda Galaxy in 1923 and 1924. Their distance established spiral nebulae well beyond the edge of the Milky Way.
Subsequent modelling of the universe explored the possibility that the cosmological constant, introduced by Einstein in his 1917 paper, may result in an expanding universe, depending on its value. Thus the Big Bang model was proposed by the Belgian priest Georges Lemaître in 1927 which was subsequently corroborated by Edwin Hubble's discovery of the redshift in 1929 and later by the discovery of the cosmic microwave background radiation by Arno Penzias and Robert Woodrow Wilson in 1964. These findings were a first step to rule out some of many alternative cosmologies.
Since around 1990, several dramatic advances in observational cosmology have transformed cosmology from a largely speculative science into a predictive science with precise agreement between theory and observation. These advances include observations of the microwave background from the COBE, WMAP and Planck satellites, large new galaxy redshift surveys including 2dfGRS and SDSS, and observations of distant supernovae and gravitational lensing. These observations matched the predictions of the cosmic inflation theory, a modified Big Bang theory, and the specific version known as the Lambda-CDM model. This has led many to refer to modern times as the "golden age of cosmology".
In 2014, the BICEP2 collaboration claimed that they had detected the imprint of gravitational waves in the cosmic microwave background. However, this result was later found to be spurious: the supposed evidence of gravitational waves was in fact due to interstellar dust.
On 1 December 2014, at the Planck 2014 meeting in Ferrara, Italy, astronomers reported that the universe is 13.8 billion years old and composed of 4.9% atomic matter, 26.6% dark matter and 68.5% dark energy.
Religious or mythological cosmology
Religious or mythological cosmology is a body of beliefs based on mythological, religious, and esoteric literature and traditions of creation and eschatology. Creation myths are found in most religions, and are typically split into five different classifications, based on a system created by Mircea Eliade and his colleague Charles Long.
Types of Creation Myths based on similar motifs:
Creation ex nihilo in which the creation is through the thought, word, dream or bodily secretions of a divine being.
Earth diver creation in which a diver, usually a bird or amphibian sent by a creator, plunges to the seabed through a primordial ocean to bring up sand or mud which develops into a terrestrial world.
Emergence myths in which progenitors pass through a series of worlds and metamorphoses until reaching the present world.
Creation by the dismemberment of a primordial being.
Creation by the splitting or ordering of a primordial unity such as the cracking of a cosmic egg or a bringing order from chaos.
Philosophy
Cosmology deals with the world as the totality of space, time and all phenomena. Historically, it has had quite a broad scope, and in many cases was found in religion. Some questions about the Universe are beyond the scope of scientific inquiry but may still be interrogated through appeals to other philosophical approaches like dialectics. Some questions that are included in extra-scientific endeavors may include: Charles Kahn, an important historian of philosophy, attributed the origins of ancient Greek cosmology to Anaximander.
What is the origin of the universe? What is its first cause (if any)? Is its existence necessary? (see monism, pantheism, emanationism and creationism)
What are the ultimate material components of the universe? (see mechanism, dynamism, hylomorphism, atomism)
What is the ultimate reason (if any) for the existence of the universe? Does the cosmos have a purpose? (see teleology)
Does the existence of consciousness have a role in the existence of reality? How do we know what we know about the totality of the cosmos? Does cosmological reasoning reveal metaphysical truths? (see epistemology)
Historical cosmologies
Table notes: the term "static" simply means not expanding and not contracting. Symbol G represents Newton's gravitational constant; Λ (Lambda) is the cosmological constant.
| Physical sciences | Astronomy basics | Astronomy |
1865464 | https://en.wikipedia.org/wiki/Carya%20ovata | Carya ovata | Carya ovata, the shagbark hickory, is a common hickory native to eastern North America, with two varieties. The trees can grow to quite a large size but are unreliable in their fruit output. The nut is consumed by wildlife and historically by Native Americans, who also used the wood.
Name
The word hickory is an aphetic form from earlier pohickory, short for even earlier pokahickory, borrowed from the Virginia Algonquian word pawcohiccora, hickory-nut meat or a nut milk drink made from it. Other names for this tree are Carolina Hickory, Scalybark Hickory, Upland Hickory, and Shellbark Hickory, with older binomial names of Carya ovata var. fraxinifolia, Carya ovata var. nuttallii, Carya ovata var. pubescens, Hicoria alba, Hicoria borealis, and Hicoria ovata.
Description
It is a large, deciduous tree, growing well over tall, and can live more than 350 years. The tallest measured shagbark, located in Savage Gulf, Tennessee, is over tall. Mature shagbarks are easy to recognize because, as their name implies, they have shaggy bark. This characteristic is, however, only found on mature trees; young specimens have smooth bark.
The leaves are long, pinnate, with five (rarely three or seven) leaflets, the terminal three leaflets much larger than the basal pair. The shagbark hickory is monoecious. Staminate flowers are borne on long-stalked catkins at the tip of old wood or in the axils of the previous season's leaves. Pistillate flowers occur in short terminal spikes. The fruit is a drupe long, an edible nut with a hard, bony shell, contained in a thick, green four-sectioned husk which turns dark and splits off at maturity in the fall. The terminal buds on the shagbark hickory are large and covered with loose scales.
C. ovata begins producing seeds at about 10 years of age, but large quantities are not produced until 40 years and will continue for at least 100. Nut production is erratic, with good crops every 3 to 5 years, in between which few or none appear and the entire crop may be lost to animal predation.
Taxonomy
The two varieties are:
Carya ovata var. ovata (northern shagbark hickory) has its largest leaflets over long and nuts long.
Carya ovata var. australis (southern shagbark hickory or Carolina hickory) has its largest leaflets under long and nuts long.
Some sources regard southern shagbark hickory as the separate species Carya carolinae-septentrionalis.
Distribution
Shagbark hickory is found throughout most of the eastern United States, but it is largely absent from the southeastern and Gulf coastal plains and lower Mississippi Delta areas. An isolated population grows in eastern Canada as far north as Lavant Township, Canadian zone 4b. Scattered locations of shagbark hickory occur in the Sierra Madre Oriental of eastern Mexico.
Shagbark hickory was introduced in Europe in the 17th century. It can still be found in Central Europe as a non-native species.
Ecology
Red squirrels, gray squirrels, raccoons, chipmunks, and mice are consumers of hickory nuts. Other consumers include black bears, gray and red foxes, rabbits, and bird species such as mallards, wood ducks, bobwhites, and wild turkey.
Uses
The nuts are edible with an excellent flavor. They are unsuitable for commercial or orchard production due to the long time it takes for a tree to produce sizable crops and unpredictable output from year to year. The nuts can be used as a substitute for the pecan in colder climates and have nearly the same culinary function. The bark is also used to flavor a maple-style syrup.
Shagbark hickory nuts were an important staple of indigenous diet. Excavation of an ancient (ca. 4350–4050 cal BP) site at Victor Mills in Columbia County, Georgia found hickory nuts, processing tools and other artifacts indicating large-scale processing and storage of nuts. Native Americans used the kernel milk to make corn cakes, kanuchi and hominy. The nuts were a significant food source for the Algonquins. Shagbark hickory wood is used for smoking meat and for making the bows of Native Americans of the northern area. The lumber is heavy, hard, and tough, weighing 63 lb/ cu ft when air-dried, and has been employed for implements and tools that require strength. These include axles, axe handles, ploughs, skis, and drum sticks.
In culture
Andrew Jackson, the seventh president of the United States, was popularly nicknamed Old Hickory, a play on the toughness of hickory wood. In 1830, he began planning the construction of his tomb at The Hermitage, his plantation in Tennessee. The grave site was surrounded by a variety of trees, including six shagbark hickories. They stood there for 168 years until a storm in 1998 demolished over 1,200 trees at the site. Work on replanting them remains an ongoing project. In modern times, shagbark hickory is rarely used as an ornamental due to its large size, slow growth, difficulty of transplanting (all Juglandaceae species have large taproots) and nut litter.
Genetics
Shagbark hickory hybridizes with pecan, Carya illinoensis, and shellbark hickory, C. laciniosa (C. x dunbarii Sarg.). Shagbark hickory has 32 chromosomes. In general, species within the genus with the same chromosome number are able to cross. Numerous hybrids among the Carya species with 32 chromosomes (pecan, bitternut, shellbark, and shagbark) have been described, though most are unproductive or have other flaws. A few hican varieties are commercially propagated.
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| Biology and health sciences | Fagales | Plants |
18579693 | https://en.wikipedia.org/wiki/A%20value | A value | A-values are numerical values used in the determination of the most stable orientation of atoms in a molecule (conformational analysis), as well as a general representation of steric bulk. A-values are derived from energy measurements of the different cyclohexane conformations of a monosubstituted cyclohexane chemical.
Substituents on a cyclohexane ring prefer to reside in the equatorial position to the axial. The difference in Gibbs free energy (ΔG) between the higher energy conformation (axial substitution) and the lower energy conformation (equatorial substitution) is the A-value for that particular substituent.
Utility
A-values help predict the conformation of cyclohexane rings. The most stable conformation will be the one which has the substituent or substituents equatorial. When multiple substituents are taken into consideration, the conformation where the substituent with the largest A-value is equatorial is favored.
The utility of A-values can be generalized for use outside of cyclohexane conformations. A-values can help predict the steric effect of a substituent. In general, the larger a substituent's A-value, the larger the steric effect of that substituent. A methyl group has an A-value of 1.74 while tert-butyl group has an A-value of ~5. Because the A-value of tert-butyl is higher, tert-butyl has a larger steric effect than methyl. This difference in steric effects can be used to help predict reactivity in chemical reactions.
Free energy considerations
Steric effects play a major role in the assignment of configurations in cyclohexanes. One can use steric hindrances to determine the propensity of a substituent to reside in the axial or equatorial plane. It is known that axial bonds are more hindered than the corresponding equatorial bonds. This is because substituents in the axial position are relatively close to two other axial substituents. This makes it very crowded when bulky substituents are oriented in the axial position. These types of steric interactions are commonly known as 1,3 diaxial interactions. These types of interactions are not present with substituents at the equatorial position.
There are generally considered three principle contributions to the conformational free energy:
Baeyer strain, defined as the strain arising from deformation of bond angles.
Pitzer strain, defined as the torsional strain arising from 1,2 interactions between groups attached to contiguous carbons,
Van der Waals interactions, which are similar to 1,3 diaxial interactions.
Enthalpic components
When comparing relative stability, 6- and 7-atom interactions can be used to approximate differences in enthalpy between conformations. Each 6-atom interaction is worth and each 7-atom interaction is worth .
Entropic components
Entropy also plays a role in a substituent's preference for the equatorial position. The entropic component is determined by the following formula:
Where σ is equal to the number of microstates available for each conformation.
Due to the larger number of possible conformations of ethyl cyclohexane, the A value is reduced from what would be predicted based purely on enthalpic terms. Due to these favorable entropic conditions, the steric relevance of an ethyl group is similar to that of a methyl substituent.
Table of A-values
Applications
Predicting reactivity
One of the original experiments performed by Winston and Holness was measuring the rate of oxidation in trans and cis substituted rings using a chromium catalyst. The large tert-butyl group used locks the conformation of each molecule, placing it equatorial (cis compound shown).
It was observed that the cis compound underwent oxidation at a much faster rate than the trans compound. The proposition was that the large hydroxyl group in the axial position was disfavored and formed the carbonyl more readily to relieve this strain. The trans compound had rates identical to those found in the monosubstituted cyclohexanol.
Approximating intramolecular force strength using A-values
Using the A-values of the hydroxyl and isopropyl subunit, the energetic value of a favorable intramolecular hydrogen bond can be calculated.
Limitations
A-Values are measured using a mono-substituted cyclohexane ring, and are an indication of only the sterics a particular substituent imparts on the molecule. This leads to a problem when there are possible stabilizing electronic factors in a different system. The carboxylic acid substituent shown below is axial in the ground state, despite a positive A-value. From this observation, it is clear that there are other possible electronic interactions that stabilize the axial conformation.
Other considerations
A-values do not predict the physical size of a molecule, only the steric effect. For example, the tert-butyl group (A-value=4.9) has a larger A-value than the trimethylsilyl group (A-value=2.5), yet the tert-butyl group actually occupies less space. This difference can be attributed to the longer length of the carbon–silicon bond as compared to the carbon–carbon bond of the tert-butyl group. The longer bond allows for less interactions with neighboring substituents, which effectively makes the trimethylsilyl group less sterically hindering, thus, lowering its A-value. This can also be seen when comparing the halogens. Bromine, iodine, and chlorine all have similar A-values even though their atomic radii differ. A-values then, predict the apparent size of a substituent, and the relative apparent sizes determine the differences in steric effects between compounds. Thus, A-values are useful tools in determining compound reactivity in chemical reactions.
| Physical sciences | Stereochemistry | Chemistry |
18580879 | https://en.wikipedia.org/wiki/Transport | Transport | Transport (in British English) or transportation (in American English) is the intentional movement of humans, animals, and goods from one location to another. Modes of transport include air, land (rail and road), water, cable, pipelines, and space. The field can be divided into infrastructure, vehicles, and operations. Transport enables human trade, which is essential for the development of civilizations.
Transport infrastructure consists of both fixed installations, including roads, railways, airways, waterways, canals, and pipelines, and terminals such as airports, railway stations, bus stations, warehouses, trucking terminals, refueling depots (including fuel docks and fuel stations), and seaports. Terminals may be used both for the interchange of passengers and cargo and for maintenance.
Means of transport are any of the different kinds of transport facilities used to carry people or cargo. They may include vehicles, riding animals, and pack animals. Vehicles may include wagons, automobiles, bicycles, buses, trains, trucks, helicopters, watercraft, spacecraft, and aircraft.
Modes
A mode of transport is a solution that makes use of a certain type of vehicle, infrastructure, and operation. The transport of a person or of cargo may involve one mode or several of the modes, with the latter case being called inter-modal or multi-modal transport. Each mode has its own advantages and disadvantages, and will be chosen on the basis of cost, capability, and route.
Governments deal with the way the vehicles are operated, and the procedures set for this purpose, including financing, legalities, and policies. In the transport industry, operations and ownership of infrastructure can be either public or private, depending on the country and mode.
Passenger transport may be public, where operators provide scheduled services, or private. Freight transport has become focused on containerization, although bulk transport is used for large volumes of durable items. Transport plays an important part in economic growth and globalization, but most types cause air pollution and use large amounts of land. While it is heavily subsidized by governments, good planning of transport is essential to make traffic flow and restrain urban sprawl.
Human-powered
Human-powered transport, a form of sustainable transport, is the transport of people or goods using human muscle-power, in the form of walking, running, and swimming. Modern technology has allowed machines to enhance human power. Human-powered transport remains popular for reasons of cost-saving, leisure, physical exercise, and environmentalism; it is sometimes the only type available, especially in underdeveloped or inaccessible regions.
Although humans are able to walk without infrastructure, the transport can be enhanced through the use of roads, especially when using the human power with vehicles, such as bicycles and inline skates. Human-powered vehicles have also been developed for difficult environments, such as snow and water, by watercraft rowing and skiing; even the air can be entered with human-powered aircraft.
Animal-powered
Animal-powered transport is the use of working animals for the movement of people and commodities. Humans may ride some of the animals directly, use them as pack animals for carrying goods, or harness them, alone or in teams, to pull sleds or wheeled vehicles.
Air
A fixed-wing aircraft, commonly called an airplane, is a heavier-than-air craft where movement of the air in relation to the wings is used to generate lift. The term is used to distinguish this from rotary-wing aircraft, where the movement of the lift surfaces relative to the air generates lift. A gyroplane is both fixed-wing and rotary wing. Fixed-wing aircraft range from small trainers and recreational aircraft to large airliners and military cargo aircraft.
Two things necessary for aircraft are air flow over the wings for lift and an area for landing. The majority of aircraft also need an airport with the infrastructure for maintenance, restocking, and refueling and for the loading and unloading of crew, cargo, and passengers. While the vast majority of aircraft land and take off on land, some are capable of take-off and landing on ice, snow, and calm water.
The aircraft is the second fastest method of transport, after the rocket. Commercial jets can reach up to , single-engine aircraft . Aviation is able to quickly transport people and limited amounts of cargo over longer distances, but incurs high costs and energy use; for short distances or in inaccessible places, helicopters can be used. As of April 28, 2009, The Guardian article notes that "the WHO estimates that up to 500,000 people are on planes at any time."
Land
Land transport covers all land-based transport systems that provide for the movement of people, goods, and services. Land transport plays a vital role in linking communities to each other. Land transport is a key factor in urban planning. It consists of two kinds, rail and road.
Rail
Rail transport is where a train runs along a set of two parallel steel rails, known as a railway or railroad. The rails are anchored perpendicular to ties (or sleepers) of timber, concrete, or steel, to maintain a consistent distance apart, or gauge. The rails and perpendicular beams are placed on a foundation made of concrete or compressed earth and gravel in a bed of ballast. Alternative methods include monorail and maglev.
A train consists of one or more connected vehicles that operate on the rails. Propulsion is commonly provided by a locomotive, that hauls a series of unpowered cars, that can carry passengers or freight. The locomotive can be powered by steam, by diesel, or by electricity supplied by trackside systems. Alternatively, some or all the cars can be powered, known as a multiple unit. Also, a train can be powered by horses, cables, gravity, pneumatics, and gas turbines. Railed vehicles move with much less friction than rubber tires on paved roads, making trains more energy efficient, though not as efficient as ships.
Intercity trains are long-haul services connecting cities; modern high-speed rail is capable of speeds up to , but this requires specially built track. Regional and commuter trains feed cities from suburbs and surrounding areas, while intra-urban transport is performed by high-capacity tramways and rapid transits, often making up the backbone of a city's public transport. Freight trains traditionally used box cars, requiring manual loading and unloading of the cargo. Since the 1960s, container trains have become the dominant solution for general freight, while large quantities of bulk are transported by dedicated trains.
Road
A road is an identifiable route, way, or path between two or more places. Roads are typically smoothed, paved, or otherwise prepared to allow easy travel; though they need not be, and historically many roads were simply recognizable routes without any formal construction or maintenance. In urban areas, roads may pass through a city or village and be named as streets, serving a dual function as urban space easement and route.
The most common road vehicle is the automobile; a wheeled passenger vehicle that carries its own motor. Other users of roads include buses, trucks, motorcycles, bicycles, and pedestrians. As of 2010, there were 1.015 billion automobiles worldwide.
Road transport offers complete freedom to road users to transfer the vehicle from one lane to the other and from one road to another according to the need and convenience. This flexibility of changes in location, direction, speed, and timings of travel is not available to other modes of transport. It is possible to provide door-to-door service only by road transport.
Automobiles provide high flexibility with low capacity, but require high energy and area use, and are the main source of harmful noise and air pollution in cities; buses allow for more efficient travel at the cost of reduced flexibility. Road transport by truck is often the initial and final stage of freight transport.
Water
Water transport is movement by means of a watercraft—such as a barge, boat, ship, or sailboat—over a body of water, such as a sea, ocean, lake, canal, or river. The need for buoyancy is common to watercraft, making the hull a dominant aspect of its construction, maintenance, and appearance.
In the 19th century, the first steam ships were developed, using a steam engine to drive a paddle wheel or propeller to move the ship. The steam was produced in a boiler using wood or coal and fed through a steam external combustion engine. Now most ships have an internal combustion engine using a slightly refined type of petroleum called bunker fuel. Some ships, such as submarines, use nuclear power to produce the steam. Recreational or educational craft still use wind power, while some smaller craft use internal combustion engines to drive one or more propellers or, in the case of jet boats, an inboard water jet. In shallow draft areas, hovercraft are propelled by large pusher-prop fans. (See Marine propulsion.)
Although it is slow compared to other transport, modern sea transport is a highly efficient method of transporting large quantities of goods. Commercial vessels, nearly 35,000 in number, carried 7.4 billion tons of cargo in 2007. Transport by water is significantly less costly than air transport for transcontinental shipping; short sea shipping and ferries remain viable in coastal areas.
Other modes
Pipeline transport sends goods through a pipe; most commonly liquid and gases are sent, but pneumatic tubes can also send solid capsules using compressed air. For liquids/gases, any chemically stable liquid or gas can be sent through a pipeline. Short-distance systems exist for sewage, slurry, water, and beer, while long-distance networks are used for petroleum and natural gas.
Cable transport is a broad mode where vehicles are pulled by cables instead of an internal power source. It is most commonly used at steep gradient. Typical solutions include aerial tramways, elevators, and ski lifts; some of these are also categorized as conveyor transport.
Spaceflight is transport outside Earth's atmosphere by means of a spacecraft. It is most frequently used for satellites placed in Earth orbit. However, human spaceflight mission have landed on the Moon and are occasionally used to rotate crew-members to space stations. Uncrewed spacecraft have also been sent to all the planets of the Solar System.
Suborbital spaceflight is the fastest of the existing and planned transport systems from a place on Earth to a distant "other place" on Earth. Faster transport could be achieved through part of a low Earth orbit or by following that trajectory even faster, using the propulsion of the rocket to steer it.
Elements
Infrastructure
Infrastructure is the fixed installations that allow a vehicle to operate. It consists of a roadway, a terminal, and facilities for parking and maintenance. For rail, pipeline, road, and cable transport, the entire way the vehicle travels must be constructed. Air and watercraft are able to avoid this, since the airway and seaway do not need to be constructed. However, they require fixed infrastructure at terminals.
Terminals such as airports, ports, and stations, are locations where passengers and freight can be transferred from one vehicle or mode to another. For passenger transport, terminals are integrating different modes to allow riders, who are interchanging between modes, to take advantage of each mode's benefits. For instance, airport rail links connect airports to the city centres and suburbs. The terminals for automobiles are parking lots, while buses and coaches can operate from simple stops. For freight, terminals act as transshipment points, though some cargo is transported directly from the point of production to the point of use.
The financing of infrastructure can either be public or private. Transport is often a natural monopoly and a necessity for the public; roads, and in some countries railways and airports, are funded through taxation. New infrastructure projects can have high costs and are often financed through debt. Many infrastructure owners, therefore, impose usage fees, such as landing fees at airports or toll plazas on roads. Independent of this, authorities may impose taxes on the purchase or use of vehicles. Because of poor forecasting and overestimation of passenger numbers by planners, there is frequently a benefits shortfall for transport infrastructure projects.
Means of transport
Animals
Animals used in transportation include pack animals and riding animals.
Vehicles
A vehicle is a non-living device that is used to move people and goods. Unlike the infrastructure, the vehicle moves along with the cargo and riders. Unless being pulled/pushed by a cable or muscle-power, the vehicle must provide its own propulsion; this is most commonly done through a steam engine, combustion engine, electric motor, jet engine, or rocket, though other means of propulsion also exist. Vehicles also need a system of converting the energy into movement; this is most commonly done through wheels, propellers, and pressure.
Vehicles are most commonly staffed by a driver. However, some systems, such as people movers and some rapid transits, are fully automated. For passenger transport, the vehicle must have a compartment, seat, or platform for the passengers. Simple vehicles, such as automobiles, bicycles, or simple aircraft, may have one of the passengers as a driver. Recently, the progress related to the Fourth Industrial Revolution has brought a lot of new emerging technologies for transportation and automotive fields such as Connected Vehicles and Autonomous Driving. These innovations are said to form future mobility, but concerns remain on safety and cybersecurity, particularly concerning connected and autonomous mobility.
Operation
Private transport is only subject to the owner of the vehicle, who operates the vehicle themselves. For public transport and freight transport, operations are done through private enterprise or by governments. The infrastructure and vehicles may be owned and operated by the same company, or they may be operated by different entities. Traditionally, many countries have had a national airline and national railway. Since the 1980s, many of these have been privatized. International shipping remains a highly competitive industry with little regulation, but ports can be public-owned.
Policy
As the population of the world increases, cities grow in size and population—according to the United Nations, 55% of the world's population live in cities, and by 2050 this number is expected to rise to 68%. Public transport policy must evolve to meet the changing priorities of the urban world. The institution of policy enforces order in transport, which is by nature chaotic as people attempt to travel from one place to another as fast as possible. This policy helps to reduce accidents and save lives.
Functions
Relocation of travelers and cargo are the most common uses of transport. However, other uses exist, such as the strategic and tactical relocation of armed forces during warfare, or the civilian mobility construction or emergency equipment.
Passenger
Passenger transport, or travel, is divided into public and private transport. Public transport is scheduled services on fixed routes, while private is vehicles that provide ad hoc services at the riders desire. The latter offers better flexibility, but has lower capacity and a higher environmental impact. Travel may be as part of daily commuting or for business, leisure, or migration.
Short-haul transport is dominated by the automobile and mass transit. The latter consists of buses in rural and small cities, supplemented with commuter rail, trams, and rapid transit in larger cities. Long-haul transport involves the use of the automobile, trains, coaches, and aircraft, the last of which have become predominantly used for the longest, including intercontinental, travel. Intermodal passenger transport is where a journey is performed through the use of several modes of transport; since all human transport normally starts and ends with walking, all passenger transport can be considered intermodal. Public transport may also involve the intermediate change of vehicle, within or across modes, at a transport hub, such as a bus or railway station.
Taxis and buses can be found on both ends of the public transport spectrum. Buses are the cheapest mode of transport but are not necessarily flexible, and taxis are very flexible but more expensive. In the middle is demand-responsive transport, offering flexibility whilst remaining affordable.
International travel may be restricted for some individuals due to legislation and visa requirements.
Medical
An ambulance is a vehicle used to transport people from or between places of treatment, and in some instances will also provide out-of-hospital medical care to the patient. The word is often associated with road-going "emergency ambulances", which form part of emergency medical services, administering emergency care to those with acute medical problems.
Air medical services is a comprehensive term covering the use of air transport to move patients to and from healthcare facilities and accident scenes. Personnel provide comprehensive prehospital and emergency and critical care to all types of patients during aeromedical evacuation or rescue operations, aboard helicopters, propeller aircraft, or jet aircraft.
Freight
Freight transport, or shipping, is a key in the value chain in manufacturing. With increased specialization and globalization, production is being located further away from consumption, rapidly increasing the demand for transport. Transport creates place utility by moving the goods from the place of production to the place of consumption. While all modes of transport are used for cargo transport, there is high differentiation between the nature of the cargo transport, in which mode is chosen. Logistics refers to the entire process of transferring products from producer to consumer, including storage, transport, transshipment, warehousing, material-handling, and packaging, with associated exchange of information. Incoterm deals with the handling of payment and responsibility of risk during transport.
Containerization, with the standardization of ISO containers on all vehicles and at all ports, has revolutionized international and domestic trade, offering a huge reduction in transshipment costs. Traditionally, all cargo had to be manually loaded and unloaded into the haul of any ship or car; containerization allows for automated handling and transfer between modes, and the standardized sizes allow for gains in economy of scale in vehicle operation. This has been one of the key driving factors in international trade and globalization since the 1950s.
Bulk transport is common with cargo that can be handled roughly without deterioration; typical examples are ore, coal, cereals, and petroleum. Because of the uniformity of the product, mechanical handling can allow enormous quantities to be handled quickly and efficiently. The low value of the cargo combined with high volume also means that economies of scale become essential in transport, and gigantic ships and whole trains are commonly used to transport bulk. Liquid products with sufficient volume may also be transported by pipeline.
Air freight has become more common for products of high value; while less than one percent of world transport by volume is by airline, it amounts to forty percent of the value. Time has become especially important in regards to principles such as postponement and just-in-time within the value chain, resulting in a high willingness to pay for quick delivery of key components or items of high value-to-weight ratio. In addition to mail, common items sent by air include electronics and fashion clothing.
Industry
Impact
Economic
Transport is a key necessity for specialization—allowing production and consumption of products to occur at different locations. Throughout history, transport has been a spur to expansion; better transport allows more trade and a greater spread of people. Economic growth has always been dependent on increasing the capacity and rationality of transport. But the infrastructure and operation of transport have a great impact on the land, and transport is the largest drainer of energy, making transport sustainability a major issue.
Due to the way modern cities and communities are planned and operated, a physical distinction between home and work is usually created, forcing people to transport themselves to places of work, study, or leisure, as well as to temporarily relocate for other daily activities. Passenger transport is also the essence of tourism, a major part of recreational transport. Commerce requires the transport of people to conduct business, either to allow face-to-face communication for important decisions or to move specialists from their regular place of work to sites where they are needed.
In lean thinking, transporting materials or work in process from one location to another is seen as one of the seven wastes (Japanese term: muda) which do not add value to a product.
Planning
Transport planning allows for high use and less impact regarding new infrastructure. Using models of transport forecasting, planners are able to predict future transport patterns. On the operative level, logistics allows owners of cargo to plan transport as part of the supply chain. Transport as a field is also studied through transport economics, a component for the creation of regulation policy by authorities. Transport engineering, a sub-discipline of civil engineering, must take into account trip generation, trip distribution, mode choice, and route assignment, while the operative level is handled through traffic engineering.
Because of the negative impacts incurred, transport often becomes the subject of controversy related to choice of mode, as well as increased capacity. Automotive transport can be seen as a tragedy of the commons, where the flexibility and comfort for the individual deteriorate the natural and urban environment for all. Density of development depends on mode of transport, with public transport allowing for better spatial use. Good land use keeps common activities close to people's homes and places higher-density development closer to transport lines and hubs, to minimize the need for transport. There are economies of agglomeration. Beyond transport, some land uses are more efficient when clustered. Transport facilities consume land, and in cities pavement (devoted to streets and parking) can easily exceed 20 percent of the total land use. An efficient transport system can reduce land waste.
Too much infrastructure and too much smoothing for maximum vehicle throughput mean that in many cities there is too much traffic and many—if not all—of the negative impacts that come with it. It is only in recent years that traditional practices have started to be questioned in many places; as a result of new types of analysis which bring in a much broader range of skills than those traditionally relied on—spanning such areas as environmental impact analysis, public health, sociology, and economics—the viability of the old mobility solutions is increasingly being questioned.
Environment
Transport is a major use of energy and burns most of the world's petroleum. This creates air pollution, including nitrous oxides and particulates, and is a significant contributor to global warming through emission of carbon dioxide, for which transport is the fastest-growing emission sector. By sub-sector, road transport is the largest contributor to global warming. Environmental regulations in developed countries have reduced individual vehicles' emissions; however, this has been offset by increases in the numbers of vehicles and in the use of each vehicle. Some pathways to reduce the carbon emissions of road vehicles considerably have been studied. Energy use and emissions vary largely between modes, causing environmentalists to call for a transition from air and road to rail and human-powered transport, as well as increased transport electrification and energy efficiency.
Other environmental impacts of transport systems include traffic congestion and automobile-oriented urban sprawl, which can consume natural habitat and agricultural lands. By reducing transport emissions globally, it is predicted that there will be significant positive effects on Earth's air quality, acid rain, smog, and climate change.
While electric cars are being built to cut down CO2 emission at the point of use, an approach that is becoming popular among cities worldwide is to prioritize public transport, bicycles, and pedestrian movement. Redirecting vehicle movement to create 20-minute neighbourhoods that promotes exercise while greatly reducing vehicle dependency and pollution. Some policies are levying a congestion charge to cars for travelling within congested areas during peak time.
Airplane emissions change depending on the flight distance. It takes a lot of energy to take off and land, so longer flights are more efficient per mile traveled. However, longer flights naturally use more fuel in total. Short flights produce the most per passenger mile, while long flights produce slightly less. Things get worse when planes fly high in the atmosphere. Their emissions trap much more heat than those released at ground level. This isn't just because of , but a mix of other greenhouse gases in the exhaust. City buses produce about 0.3 kg of for every mile traveled per passenger. For long-distance bus trips (over 20 miles), that pollution drops to about 0.08 kg of per passenger mile. On average, commuter trains produce around 0.17 kg of for each mile traveled per passenger. Long-distance trains are slightly higher at about 0.19 kg of per passenger mile. The fleet emission average for delivery vans, trucks and big rigs is per gallon of diesel consumed. Delivery vans and trucks average about 7.8 mpg (or 1.3 kg of per mile) while big rigs average about 5.3 mpg (or 1.92 kg of per mile).
Sustainable development
The United Nations first formally recognized the role of transport in sustainable development in the 1992 United Nations Earth summit. In the 2012 United Nations World Conference, global leaders unanimously recognized that transport and mobility are central to achieving the sustainability targets. In recent years, data has been collected to show that the transport sector contributes to a quarter of the global greenhouse gas emissions, and therefore sustainable transport has been mainstreamed across several of the 2030 Sustainable Development Goals, especially those related to food, security, health, energy, economic growth, infrastructure, and cities and human settlements. Meeting sustainable transport targets is said to be particularly important to achieving the Paris Agreement.
There are various Sustainable Development Goals (SDGs) that are promoting sustainable transport to meet the defined goals. These include SDG 3 on health (increased road safety), SDG 7 on energy, SDG 8 on decent work and economic growth, SDG 9 on resilient infrastructure, SDG 11 on sustainable cities (access to transport and expanded public transport), SDG 12 on sustainable consumption and production (ending fossil fuel subsidies), and SDG 14 on oceans, seas, and marine resources.
History
Natural
Humans' first ways to move included walking, running, and swimming. The domestication of animals introduced a new way to lay the burden of transport on more powerful creatures, allowing the hauling of heavier loads, or humans riding animals for greater speed and duration. Inventions such as the wheel and the sled (U.K. sledge) helped make animal transport more efficient through the introduction of vehicles.
The first forms of road transport involved animals, such as horses (domesticated in the 4th or the 3rd millennium BCE), oxen (from about 8000 BCE), or humans carrying goods over dirt tracks that often followed game trails.
Water transport
Water transport, including rowed and sailed vessels, dates back to time immemorial and was the only efficient way to transport large quantities or over large distances prior to the Industrial Revolution. The first watercraft were canoes cut out from tree trunks. Early water transport was accomplished with ships that were either rowed or used the wind for propulsion, or a combination of the two. The importance of water has led to most cities that grew up as sites for trading being located on rivers or on the sea-shore, often at the intersection of two bodies of water.
Mechanical
Until the Industrial Revolution, transport remained slow and costly, and production and consumption gravitated as close to each other as feasible. The Industrial Revolution in the 19th century saw several inventions fundamentally change transport. With telegraphy, communication became instant and independent of the transport of physical objects. The invention of the steam engine, closely followed by its application in rail transport, made land transport independent of human or animal muscles. Both speed and capacity increased, allowing specialization through manufacturing being located independently of natural resources. The 19th century also saw the development of the steam ship, which sped up global transport.
With the development of the combustion engine and the automobile around 1900, road transport became more competitive again, and mechanical private transport originated. The first "modern" highways were constructed during the 19th century with macadam. Later, tarmac and concrete became the dominant paving materials.
In 1903 the Wright brothers demonstrated the first successful controllable airplane, and after World War I (1914–1918) aircraft became a fast way to transport people and express goods over long distances.
After World War II (1939–1945) the automobile and airlines took higher shares of transport, reducing rail and water to freight and short-haul passenger services. Scientific spaceflight began in the 1950s, with rapid growth until the 1970s, when interest dwindled. In the 1950s the introduction of containerization gave massive efficiency gains in freight transport, fostering globalization. International air travel became much more accessible in the 1960s with the commercialization of the jet engine. Along with the growth in automobiles and motorways, rail and water transport declined in relative importance. After the introduction of the Shinkansen in Japan in 1964, high-speed rail in Asia and Europe started attracting passengers on long-haul routes away from the airlines.
Early in U.S. history, private joint-stock corporations owned most aqueducts, bridges, canals, railroads, roads, and tunnels. Most such transport infrastructure came under government control in the late 19th and early 20th centuries, culminating in the nationalization of inter-city passenger rail-service with the establishment of Amtrak. Recently, however, a movement to privatize roads and other infrastructure has gained some ground and adherents.
| Technology | Technology | null |
18581082 | https://en.wikipedia.org/wiki/Tarantula | Tarantula | Tarantulas comprise a group of large and often hairy spiders of the family Theraphosidae. , 1,100 species have been identified, with 166 genera. The term "tarantula" is usually used to describe members of the family Theraphosidae, although many other members of the same infraorder (Mygalomorphae) are commonly referred to as "tarantulas" or "false tarantulas". Some of the more common species have become popular in the exotic pet trade. Many New World species kept as pets have setae known as urticating hairs that can cause irritation to the skin, and in extreme cases, cause damage to the eyes.
Overview
Like all arthropods, the tarantula is an invertebrate that relies on an exoskeleton for muscular support. Like other Arachnida, a tarantula's body comprises two main parts, the prosoma (or cephalothorax) and the opisthosoma (or abdomen). The prosoma and opisthosoma are connected by the pedicel, or pregenital somite. This waist-like connecting piece is actually part of the prosoma and gives the opisthosoma a wide range of motion relative to the prosoma.
Depending on the species, the body length of tarantulas ranges from about with leg spans of . Leg span is determined by measuring from the tip of the back leg to the tip of the front leg on the opposite side. Some of the largest species of tarantula may weigh over ; the largest of all, the goliath birdeater (Theraphosa blondi) from Venezuela and Brazil, has been reported to attain a weight of and a leg-span up to , males being longer and females greater in girth. The fang size of this tarantula reaches a maximum of .
Theraphosa apophysis (the pinkfoot goliath) was described 187 years after the goliath birdeater, so its characteristics are not as well attested. T. blondi is generally thought to be the heaviest tarantula, and T. apophysis has the greatest leg span. Two other species, Lasiodora parahybana (the Brazilian salmon birdeater) and Lasiodora klugi, rival the size of the two goliath spiders.
Most species of North American tarantulas are brown. Elsewhere, species have been found that variously display cobalt blue (Cyriopagopus lividus), black with white stripes (Aphonopelma seemanni), yellow leg markings (Eupalaestrus campestratus), metallic blue legs with vibrant orange abdomen and green prosoma (Chromatopelma cyaneopubescens). Their natural habitats include savanna, grassland such as in the pampas, rainforest, desert, scrubland, mountains, and cloud forest. They are generally classed among the terrestrial types. They are burrowers that live in the ground.
Tarantulas are becoming increasingly popular as pets and some species are readily available in captivity.
Identification
Tarantulas can be confused with other members of the order Mygalomorphae, such as trapdoor spiders, funnel-web spiders and purseweb spiders. They can also be confused with some members of the order Araneomorphae such as the Lycosidae family. There are multiple ways to identify a tarantula. First the hairs: in the Americas most tarantulas have urticating hairs, though some, such as the Hemirrhagus genus, lack these. The hairs are usually more noticeable than with most other spiders. Another is the size, as tarantulas tend to be bigger, but this is again not a failproof way. They also do not use their webs for hunting, instead using them as building material or tripwire.
One of the most decisive ways to tell is by looking at their fangs. Tarantula fangs face downwards, as opposed to those of true spiders, which face each other, allowing them to make pincerlike motions. They also own two book lungs, as opposed to true spiders which only have one. Their lifespan is also longer than most spiders.
Etymology
The spider originally bearing the name tarantula was Lycosa tarantula, a species of wolf spider native to Mediterranean Europe. The name is derived from the southern Italian town of Taranto. The term tarantula was subsequently applied to almost any large, unfamiliar species of ground-dwelling spider, in particular to the Mygalomorphae and especially the New World Theraphosidae. Compared to tarantulas, wolf spiders are not particularly large or hairy, and so among English speakers in particular, usage eventually shifted in favour of the Theraphosidae, even though they are not closely related to wolf spiders at all, being in a different infraorder.
The name tarantula is also incorrectly applied to other large-bodied spiders, including the purseweb spiders or atypical tarantulas, the funnel-webs (Dipluridae and Hexathelidae), and the dwarf tarantulas. These spiders are related to tarantulas (all being mygalomorphs) but fall into different families from them. Huntsman spiders of the family Sparassidae have also been termed tarantulas because of their large size, when, in fact, they are not related. Instead, huntsman spiders belong to the infraorder Araneomorphae.
The element pelma in genus names
Many theraphosid genera have names, either accepted or synonymous, containing the element pelma. This can be traced back to Carl Ludwig Koch in 1850, who in describing his new genus Eurypelma wrote, "" (). German arachnologists use the word to refer to the tarsus (the last article of a spider's leg). Translations of into Latin use the word . Hence in English arachnological terminology, Koch meant 'the scopula of the base of the tarsus very wide'. Eury- is derived from the Greek (), meaning 'wide', while () means 'the sole of the foot', paralleling Koch's use of (in modern spelling). Thus Eurypelma literally means 'wide footsole'; however, arachnologists have conventionally taken pelma in such names to refer to the scopula, so producing the meaning 'with a wide scopula'.
Other genus names or synonyms that Estrada-Alvarez and Cameron regard as having 'footsole' or 'scopula' meanings include:
Acanthopelma – Greek () 'thorn, spine'; overall meaning 'spiny footsole'
Brachypelma – Greek () 'short'; overall meaning 'short scopula'
Metriopelma – Greek () 'of moderate size'; overall meaning 'medium length scopula'
Schizopelma – from the Greek origin combining form schizo- () 'split'; overall meaning 'split footsole'
Sericopelma – Greek () 'silky'; overall meaning 'silken scopula'
Later, particularly following genus names published by R.I. Pocock in 1901, the element pelma appears to have become synonymous with 'theraphosid'. For example, the author of Cardiopelma writes, "" ('Cardiopelma refers to the female genitalia that evoke the shape of a heart'), with no reference to either 'footsole' or 'scopula'. Names interpreted in this way include:
Aphonopelma – Greek () 'soundless'; overall meaning 'theraphosid without sound'
Cardiopelma – Greek () 'heart'; overall meaning 'heart theraphosid' (referring to the heart-shaped female genitalia)
Clavopelma – Latin 'club'; overall meaning 'theraphosid with club-shaped hairs'
Delopelma – Greek () 'clear, obvious, visible, conspicuous, plain'; overall meaning 'theraphosid without plumose hair'
Gosipelma – the element gosi- means 'desert', relating to the Gosiute people; overall meaning 'desert theraphosid'
Spelopelma – Greek () 'cave'; overall meaning 'cave theraphosid'
Distribution
Tarantulas of various species occur throughout the United States, Mexico, in Central America, and throughout South America. Other species occur variously throughout Africa, much of Asia (including the Ryukyu Islands in southern Japan), and all of Australia. In Europe, some species occur in Spain, Portugal, Turkey, southern Italy, and Cyprus.
Habits
Some genera of tarantulas hunt prey primarily in trees; others hunt on or near the ground. All tarantulas can produce silk; while arboreal species typically reside in a silken "tube tent", terrestrial species line their burrows with silk to stabilize the burrow wall and facilitate climbing up and down. Tarantulas mainly eat large insects and other arthropods such as centipedes, millipedes, and other spiders, using ambush as their primary method of prey capture. Armed with their massive, powerful chelicerae tipped with long, chitinous fangs, tarantulas are well-adapted to killing other large arthropods. The biggest tarantulas sometimes kill and consume small vertebrates such as lizards, mice, bats, birds, and small snakes.
Appendages
The eight legs, the two chelicerae with their fangs, and the pedipalps are attached to the prosoma. The chelicerae are two double-segmented appendages located just below the eyes and directly forward of the mouth. The chelicerae contain the venom glands that vent through the fangs. The fangs are hollow extensions of the chelicerae that inject venom into prey or animals that the tarantula bites in defense, and they are also used to masticate. These fangs are articulated so that they can extend downward and outward in preparation to bite or can fold back toward the chelicerae as a pocket knife blade folds back into its handle. The chelicerae of a tarantula completely contain the venom glands and the muscles that surround them, and can cause the venom to be forcefully injected into prey.
The pedipalpi are two six-segmented appendages connected to the prosoma near the mouth and protruding on either side of both chelicerae. In most species of tarantulas, the pedipalpi contain sharp, jagged plates used to cut and crush food often called the coxae or maxillae. As with other spiders, the terminal portions of the pedipalpi of males function as part of their reproductive system. Male spiders spin a silken platform (sperm web) on the ground onto which they release semen from glands in their opisthosoma. Then they insert their pedipalps into the semen, absorb the semen into the pedipalps, and later insert the pedipalps (one at a time) into the reproductive organ of the female, which is located in her abdomen. The terminal segments of the pedipalps of male tarantulas are moderately larger in circumference than those of a female tarantula. Male tarantulas have special spinnerets surrounding the genital opening. Silk for the sperm web of the tarantula is exuded from these special spinnerets.
A tarantula has four pairs of legs and two additional pairs of appendages. Each leg has seven segments, which from the prosoma out are: coxa, trochanter, femur, patella, tibia, tarsus and pretarsus, and claw. Two or three retractable claws at the end of each leg are used to grip surfaces for climbing. Also on the end of each leg, surrounding the claws, is a group of bristles, called the scopula, which help the tarantula to grip better when climbing surfaces such as glass. The fifth pair is the pedipalps, which aid in feeling, gripping prey, and mating in the case of a mature male. The sixth pair of appendages is the chelicerae and their attached fangs. When walking, a tarantula's first and third legs on one side move at the same time as the second and fourth legs on the other side of its body. The muscles in a tarantula's legs cause the legs to bend at the joints, but to extend a leg, the tarantula increases the pressure of haemolymph entering the leg.
Tarantulas, like almost all other spiders, have their primary spinnerets at the end of the opisthosoma. Unlike most spider species in the infraorder Araneomorphae, which includes the majority of extant spider species, and most of which have six, tarantula species have two or four spinnerets. Spinnerets are flexible, tube-like structures from which the spider exudes its silk. The tip of each spinneret is called the spinning field. Each spinning field is covered by as many as 100 spinning tubes through which silk is exuded. As the silk is pulled out of the spinnerets, the shear forces cause proteins in the silk to crystallize, transforming it from a liquid to a solid thread.
Digestive system
The tarantula's mouth is located under its chelicerae on the lower front part of its prosoma. The mouth is a short, straw-shaped opening that can only suck, meaning that anything taken into it must be in liquid form. Prey with large amounts of solid parts, such as mice, must be crushed and ground up or predigested, which is accomplished by coating the prey with digestive juices secreted from openings in the chelicerae.
The tarantula's digestive organ (stomach) is a tube that runs the length of its body. In the prosoma, this tube is wider and forms the sucking stomach. When the sucking stomach's powerful muscles contract, the stomach is increased in cross-section, creating a strong sucking action that permits the tarantula to suck its liquefied prey up through the mouth and into the intestines. Once the liquefied food enters the intestines, it is broken down into particles small enough to pass through the intestine walls into the hemolymph (blood stream), where it is distributed throughout the body. After feeding, the leftovers are formed into a small ball by the tarantula and thrown away. In a terrarium, they often put them into the same corner.
Nervous system
A tarantula's central nervous system (brain) is located in the bottom of the inner prosoma. A tarantula perceives its surroundings primarily via sensory organs called setae (bristles or spines, sometimes referred to as hairs). Although a tarantula has eight eyes like most spiders, touch is its keenest sense, and in hunting, it primarily depends on vibrations given off by the movements of its prey. A tarantula's setae are very sensitive organs and are used to sense chemical signatures, vibrations, wind direction, and possibly even sound. Tarantulas are also very responsive to the presence of certain chemicals such as pheromones.
The eyes are located above the chelicerae on the forward part of the prosoma. They are small and usually set in two rows of four. Most tarantulas are not able to see much more than light, darkness, and motion. Arboreal tarantulas generally have better vision compared with terrestrial tarantulas.
Respiratory system
All types of tarantulas have two sets of book lungs (breathing organs); the first pair is located in a cavity inside the lower front part of the abdomen near where the abdomen connects to the cephalothorax, and the second pair is slightly farther back on the abdomen. Air enters the cavity through a tiny slit on each side of and near the front of the abdomen. Each lung consists of 15 or more thin sheets of folded tissue arranged like the pages of a book. These sheets of tissue are supplied by blood vessels. As air enters each lung, oxygen is taken into the blood stream through the blood vessels in the lungs. Needed moisture may also be absorbed from humid air by these organs.
Circulatory system
A tarantula's blood is unique (not only in appearance); an oxygen-transporting protein is present (the copper-based hemocyanin), but not enclosed in blood cells such as the erythrocytes of mammals. A tarantula's blood is not true blood, but rather a liquid called hemolymph (or haemolymph). At least four types of hemocytes, or hemolymph cells, are known.
The tarantula's heart is a long, slender tube located along the top of the opisthosoma. The heart is neurogenic as opposed to myogenic, so nerve cells instead of muscle cells initiate and coordinate the heart. It pumps hemolymph to all parts of the body through open passages often referred to as sinuses, and not through a circular system of blood vessels. If the exoskeleton is breached, loss of hemolymph will kill the spider unless the wound is small enough that the hemolymph can dry and close it.
Predators
Despite their large size and fearsome appearance and reputation, tarantulas themselves are prey for many other animals. The most specialized of these predators are large members of the wasp family Pompilidae such as the wasp Hemipepsis ustulata. These wasps are called "tarantula hawks". The largest tarantula hawks, such as those in the genus Pepsis, track, attack, and kill large tarantulas. They use olfaction to find the lair of a tarantula. The wasp must deliver a sting to the underside of the spider's cephalothorax, exploiting the thin membrane between the basal leg segments. This paralyzes the spider, and the wasp then drags it back into its burrow before depositing an egg on the prey's abdomen. The wasp then seals the spider in its burrow and flies off to search for more hosts. The wasp egg hatches into a larva and feeds on the spider's inessential parts, and as it approaches pupation, it consumes the remainder. Other arthropods, such as large scorpions and giant centipedes, are also known to prey on tarantulas.
Tarantulas are also preyed upon by a wide variety of vertebrates. Many of these, including lizards, frogs, birds, snakes and mammals, are generalist predators of all kinds of large arthropods. Mammals that have been known to prey on tarantulas, such as the coati, kinkajou, and opossum in the New World, and mongooses and the honey badger in the Old World, are often immune to the venom of their arthropod prey.
Humans also consume tarantulas for food in their native ranges. They are considered a delicacy in certain cultures (e.g. Venezuela and Cambodia). They can be roasted over an open fire to remove the bristles (described further below) and then eaten.
Tarantulas have evolved specialized bristles, or setae, to defend themselves against predators. Besides the normal bristles covering the body, some tarantulas also have a dense covering of irritating bristles called urticating hairs, on the opisthosoma, that they sometimes use as protection against enemies. These bristles are present on most New World species, but not on any specimens from the Old World. Urticating hairs are usually kicked off the abdomen by the tarantula, but some may simply rub the abdomen against the target, like the genus Avicularia. These fine bristles are barbed and serve to irritate. They can be lethal to small animals such as rodents. Some people are sensitive to these bristles, and develop serious itching and rashes at the site. Exposure of the eyes and respiratory system to urticating hairs should be strictly avoided. Species with urticating hairs can kick these bristles off; they are flicked into the air at a target using their back pairs of legs. Tarantulas also use these bristles for other purposes, such as to mark territory or to line their shelters (the latter such practice may discourage flies from feeding on the spiderlings). Urticating hairs do not grow back, but are replaced with each molt. The intensity, number, and flotation of the bristles depends on the species of tarantula.
To predators and other enemies, these bristles can range from being lethal to simply being a deterrent. With humans, they can cause irritation to eyes, nose, and skin, and more dangerously, the lungs and airways, if inhaled. The symptoms range from species to species, from person to person, from a burning itch to a minor rash. In some cases, tarantula bristles have caused permanent damage to human eyes.
Some setae are used to stridulate, which makes a hissing sound. These bristles are usually found on the chelicerae. Stridulation seems to be more common in Old World species.
Bites and urticating bristles
All tarantulas are venomous. Although their venom is not deadly to humans, some bites cause serious discomfort that might persist for several days. In general, the effects of the bites of all kinds of tarantula are not well known. While the bites of many species are known to be no worse than a wasp sting, accounts of bites by some species are reported to be very painful and to produce intense spasms that may recur over a period of several days; the venom of the African tarantula Pelinobius muticus also causes strong hallucinations. For Poecilotheria species, researchers have described more than 20 bites with the delayed onset of severe and diffuse muscle cramps, lasting for several days, that in most cases resolved completely with the use of benzodiazepines and magnesium. In all cases, seeking medical aid is advised. Because other proteins are included when a toxin is injected, some individuals may suffer severe symptoms due to an allergic reaction rather than to the venom. Such allergic effects can be life-threatening. Additionally, the large fangs of a tarantula can inflict painful puncture wounds, which can lead to secondary bacterial infections if not properly treated.
Before biting, a tarantula may signal its intention to attack by rearing up into a "threat posture", which may involve raising its prosoma and lifting its front legs into the air, spreading and extending its fangs, and (in certain species) making a loud hissing by stridulating. Tarantulas often hold this position for longer than the duration of the original threat. Their next step, without biting, may be to slap down on the intruder with their raised front legs. If that response fails to deter the attacker, the tarantulas of the Americas may next turn away and flick urticating hairs toward the pursuing predator. The next response may be to leave the scene entirely, but especially if no line of retreat is available, their final response may also be to whirl suddenly and bite. Some tarantulas are well known to give "dry bites", i.e., they may defensively bite some animal that intrudes on their space and threatens them, but they do not pump venom into the wound.
New-world tarantulas—those indigenous to the Americas—have bites that generally pose little threat to humans (other than causing localized pain). Most of them are equipped with urticating hairs on their abdomens, and almost always throw these barbed bristles as the first line of defense. These bristles irritate sensitive areas of the body and especially seem to target curious animals that may sniff these bristles into the mucous membranes of the nose. Some species have more effective urticating bristles than others. The goliath birdeater is known for its particularly irritating urticating bristles. They can penetrate the cornea, so eye protection should be worn when handling such tarantulas.
Old World tarantulas have no urticating bristles and are more likely to attack when disturbed. They often have more potent, medically significant venom, and are faster and much more nervous and defensive than New World species.
Some dangerous spider species are related to tarantulas and are frequently confused with them. A popular urban legend maintains that deadly varieties of tarantula exist somewhere in South America. This claim is often made without identifying a particular spider, although the "banana tarantula" is sometimes named. A likely candidate for the true identity of this spider is the dangerous Brazilian wandering spider (Phoneutria fera) of the family Ctenidae, as it is sometimes found hiding in clusters of bananas and is one of several spiders called "banana spiders". It is not technically a tarantula, but it is fairly large (4- to 5-inch legspan), somewhat ″hairy″, and is highly venomous to humans. Another dangerous type of spiders that have been confused with tarantulas are the Australian funnel-web spiders. The best known species of these is the Sydney funnel-web spider (Atrax robustus) a spider that is aggressive, highly venomous, and (prior to the development of antivenom in the 1980s) was responsible for numerous deaths in Australia. These spiders are members of the same infraorder as tarantulas, Mygalomorphae. Some Australians use the slang term "triantelope" (a corruption of the incorrect term tarantula, which is also used) for large, ″hairy″, and harmless members of the huntsman spider family, which are often found on interior household walls and in automobiles.
Sexual dimorphism
Some tarantula species exhibit pronounced sexual dimorphism. Males tend to be smaller (especially their abdomens, which can appear quite narrow) and may be dull in color when compared to their female counterparts, as in the species Haplopelma lividum. Mature male tarantulas also may have tibial hooks on their front legs, which are used to restrain the female's fangs during copulation. Males typically have longer legs than the females.
A juvenile male's sex can be determined by looking at a cast exuvia for epiandrous fusillae or spermathecae. Females possess spermathecae, except for the species Sickius longibulbi and Encyocratella olivacea. Males have much shorter lifespans than females because they die relatively soon after maturing. Few live long enough for a postultimate molt, which is unlikely in natural habitats because they are vulnerable to predation, but has happened in captivity, though rarely. Most males do not live through this molt, as they tend to get their emboli, mature male sexual organs on pedipalps, stuck in the molt. Most tarantula fanciers regard females as more desirable as pets due to their much longer lifespans. Wild-caught tarantulas are often mature males because they wander out in the open and are more likely to be caught.
Life cycle
Like other spiders, tarantulas have to shed their exoskeleton periodically as they grow, a process called molting. A young tarantula may do this several times a year as a part of the maturation process, while full-grown specimens only molt once a year or less, or sooner, to replace lost limbs or lost urticating hairs. It is visibly apparent that molting is imminent when the exoskeleton takes on a darker shade. If a tarantula previously used its urticating hairs, the bald patch turns from a peach color to deep blue. The tarantula also stops feeding and becomes more lethargic during this time.
While most Tarantulas species take between two and five years to reach sexual maturity, some species can take up to 10 years. Upon reaching adulthood, males typically have an 18-month period left to live so immediately go in search of a female mate. Although females continue to molt after reaching maturity, males rarely do again once they reach adulthood. Those that do often can become stuck during the molting process due to their sexual organs and die.
Females can live for 30 to 40 years. Grammostola rosea spiders, which eat once or twice a week, have lived up to 20 years in captivity. Some have survived on water alone for up to two years.
Reproduction
After reaching sexual maturity, a female tarantula normally mates and lays eggs once per year, although they do not always do so.
As with other spiders, the mechanics of intercourse are quite different from those of mammals. Once a male spider reaches maturity and becomes motivated to mate, he weaves a web mat on a flat surface. The spider then rubs his abdomen on the surface of this mat, and in so doing, releases a quantity of semen. He may then insert his pedipalps (short, leg-like appendages between the chelicerae and front legs) into the pool of semen. The pedipalps absorb the semen and keep it viable until a mate can be found. When a male spider detects the presence of a female, the two exchange signals to establish that they are of the same species. These signals may also lull the female into a receptive state. If the female is receptive, then the male approaches her and inserts his pedipalps into an opening in the lower surface of her abdomen, the opisthosoma. After the semen has been transferred to the receptive female's body, the male swiftly leaves the scene before the female recovers her appetite. Although females may show some aggression after mating, the male rarely becomes a meal.
Females deposit 50 to 2,000 eggs, depending on the species, in a silken egg sac and guard it for six to eight weeks. During this time, the females stay very close to the egg sacs and become more aggressive. Within most species, the females turn the egg sac often, which is called brooding. This keeps the eggs from deforming due to sitting in one position too long. The young spiderlings remain in the nest for some time after hatching, where they live off the remains of their yolk sacs before dispersing.
Taxonomy
Linnaeus placed all spiders in a single genus, Aranea. In 1802, Charles Athanase Walckenaer separated mygalomorph spiders into a separate genus, Mygale, leaving all other spiders in Aranea. However, Mygale had already been used in 1800 by Georges Cuvier for a genus of mammals (in Greek, mygale means "shrew"). Accordingly, in 1869, Tamerlan Thorell used the family name "Theraphosoidae" (modern Theraphosidae) for the mygalomorph spiders known to him, rather than "Mygalidae" (as used, for example, by John Blackwall). Thorell later split the family into a number of genera, including Theraphosa.
Subfamilies
A 2019 phylogenomic study recognized 12 subfamilies, one (Ischnocolinae) known not to be monophyletic.
Aviculariinae
Eumenophorinae
Harpactirinae
Ischnocolinae
Ornithoctoninae
Poecilotheriinae
Psalmopoeinae
Schismatothelinae
Selenocosmiinae
Stromatopelminae
Theraphosinae
Thrigmopoeinae
The relationship between the subfamilies found in the study is shown in the following cladogram. The dual placing of Ischnocolinae is highlighted.
All the species that possess urticating hairs and have been seen to use them in bombardment behavior are placed in the "bombardier clade", although not all species in the included subfamilies possess such hairs (all Schismatothelinae lack them as do most Psalmopoeinae genera). It is not clear whether the possession of urticating hairs was an ancestral trait of the clade, and has been lost in some species, or whether it represents multiple gains. Foley et al. suggested that the second hypothesis appeared to be better supported.
Other subfamilies that have been recognized include:
Acanthopelminae – may be treated as synonymous with Ischnocolinae
Selenogyrinae
Spelopelminae – typically not accepted, Hemirrhagus being treated as Theraphosinae
Genera
, the World Spider Catalog accepted the following genera:
Abdomegaphobema Sherwood, Gabriel, Peñaherrera-R., Léon-E., Cisneros-Heredia, Brescovit & Lucas, 2023
Acanthopelma F. O. Pickard-Cambridge, 1897 – Guyana
Acanthoscurria Ausserer, 1871 – South America, Guatemala
Acentropelma Pocock, 1901 – Belize, Mexico, Guatemala
Aenigmarachne Schmidt, 2005 – Costa Rica
Agnostopelma Pérez-Miles & Weinmann, 2010 – Colombia
Aguapanela Perafán & Cifuentes, 2015
Amazonius Cifuentes & Bertani, 2022 - South America
Annandaliella Hirst, 1909 – India
Anoploscelus Pocock, 1897 – Uganda, Tanzania, Rwanda
Anqasha Sherwood & Gabriel, 2022 - Peru
Antikuna Kaderka, Ferretti, West, Lüddecke & Hüsser, 2021 - Peru
Antillena Bertani, Huff & Fukushima, 2017 – Dominican Republic
Aphonopelma Pocock, 1901 – North America, Central America
Augacephalus Gallon, 2002 – South Africa, Mozambique, Eswatini
Avicularia Lamarck, 1818 – South America, Trinidad and Tobago, Panama
Bacillochilus Gallon, 2010 – Angola
Batesiella Pocock, 1903 – Cameroon
Bermejoa Gabriel, Sherwood & Pérez-Miles, 2023
Birupes Gabriel & Sherwood, 2019 – Malaysia
Bistriopelma Kaderka, 2015 – Peru
Bonnetina Vol, 2000 – Mexico
Brachionopus Pocock, 1897 – South Africa
Brachypelma Simon, 1891 – Mexico, Costa Rica, Guatemala
Bumba Pérez-Miles, Bonaldo & Miglio, 2014 – Brazil, Venezuela, Ecuador
Cardiopelma Vol, 1999 – Unknown
Caribena Fukushima & Bertani, 2017 – Cuba
Catanduba Yamamoto, Lucas & Brescovit, 2012 – Brazil
Catumiri Guadanucci, 2004 – South America
Ceratogyrus Pocock, 1897 – Africa
Chaetopelma Ausserer, 1871 – Asia, Greece, Africa
Chilobrachys Karsch, 1892 – Asia
Chinchaysuyu Ferretti, Chaparro, Ochoa & West, 2023
Chromatopelma Schmidt, 1995 – Venezuela
Cilantica Mirza, 2024 - India
Citharacanthus Pocock, 1901 – Cuba, Central America, Mexico
Citharognathus Pocock, 1895 – Indonesia
Clavopelma Chamberlin, 1940 – Mexico
Coremiocnemis Simon, 1892 – Malaysia, Indonesia, Australia
Cotztetlana Mendoza, 2012 – Mexico
Crassicrus Reichling & West, 1996 – Mexico, Belize
Cubanana Ortiz, 2008 – Cuba
Cyclosternum Ausserer, 1871 – South America, Mexico, Costa Rica
Cymbiapophysa Gabriel & Sherwood, 2020
Cyriocosmus Simon, 1903 – South America, Trinidad and Tobago
Cyriopagopus Simon, 1887 – Asia
Cyrtogrammomma Pocock, 1895 - Guyana and Brazil
Cyrtopholis Simon, 1892 – Caribbean
Davus O. Pickard-Cambridge, 1892 – Central America, Mexico
Dolichothele Mello-Leitão, 1923 – Brazil, Bolivia
Dugesiella Pocock, 1901 - Mexico
Encyocratella Strand, 1907 – Tanzania
Encyocrates Simon, 1892 – Madagascar
Ephebopus Simon, 1892 – Suriname, Brazil
Euathlus Ausserer, 1875 – Chile, Argentina
Eucratoscelus Pocock, 1898 – Kenya, Tanzania
Eumenophorus Pocock, 1897 – Sierra Leone
Eupalaestrus Pocock, 1901 – South America
Euphrictus Hirst, 1908 – Cameroon, Congo
Euthycaelus Simon, 1889 – Colombia, Venezuela
Grammostola Simon, 1892 – South America
Guyruita Guadanucci, Lucas, Indicatti & Yamamoto, 2007 – Brazil, Venezuela
Hapalopus Ausserer, 1875 – South America, Panama
Hapalotremus Simon, 1903 – Bolivia, Peru, Argentina
Haploclastus Simon, 1892 – India
Haplocosmia Schmidt & von Wirth, 1996 – Nepal
Harpactira Ausserer, 1871 – South Africa, Namibia
Harpactirella Purcell, 1902 – South Africa, Morocco
Hemirrhagus Simon, 1903 – Mexico
Heterophrictus Pocock, 1900 – India
Heteroscodra Pocock, 1900 – Cameroon, Gabon, Congo
Heterothele Karsch, 1879 – Africa, Argentina
Holothele Karsch, 1879 – Caribbean, South America
Homoeomma Ausserer, 1871 – South America
Hysterocrates Simon, 1892 – Africa
Idiothele Hewitt, 1919 – South Africa
Iridopelma Pocock, 1901 – Brazil
Ischnocolus Ausserer, 1871 – Africa, Asia, Brazil, Europe
Isiboroa Gabriel, Sherwood & Pérez-Miles, 2023
Kankuamo Perafán, Galvis & Pérez-Miles, 2016
Kochiana Fukushima, Nagahama & Bertani, 2008 – Brazil
Lampropelma Simon, 1892 – Indonesia, Malaysia, Singapore
Lasiocyano Galleti-Lima, Hamilton, Borges & Guadanucci, 2023 – Brazil
Lasiodora C. L. Koch, 1850 – South America, Costa Rica
Lasiodorides Schmidt & Bischoff, 1997 – Ecuador, Peru
Longilyra Gabriel, 2014 – El Salvador
Loxomphalia Simon, 1889 – Tanzania
Loxoptygus Simon, 1903 – Ethiopia
Lyrognathus Pocock, 1895 – Indonesia, India, Malaysia
Magnacarina Mendoza, Locht, Kaderka, Medina & Pérez-Miles, 2016 – Mexico
Mascaraneus Gallon, 2005 – Mauritius
Megaphobema Pocock, 1901 – Costa Rica, Colombia, Ecuador
Melloina Brignoli, 1985 – Panama, Venezuela
Melognathus Chamberlin, 1917
Metriopelma Becker, 1878 – Mexico
Miaschistopus Pocock, 1897 – Venezuela
Monocentropus Pocock, 1897 – Yemen, Madagascar
Munduruku Miglio, Bonaldo & Pérez-Miles, 2013
Murphyarachne Sherwood & Gabriel, 2022 - Peru
Mygalarachne Ausserer, 1871 – Honduras
Myostola Simon, 1903 – Gabon, Cameroon
Neischnocolus Petrunkevitch, 1925 – Panama
Neoheterophrictus Siliwal & Raven, 2012 – India
Neoholothele Guadanucci & Weinmann, 2015 – Colombia, Trinidad and Tobago, Venezuela
Neostenotarsus Pribik & Weinmann, 2004 – French Guiana
Nesiergus Simon, 1903 – Seychelles
Nesipelma Schmidt & Kovařík, 1996 – St. Kitts and Nevis
Nhandu Lucas, 1983 – Brazil, Paraguay
Omothymus Thorell, 1891 – Malaysia
Ornithoctonus Pocock, 1892 – Myanmar, Thailand
Orphnaecus Simon, 1892 – Papua New Guinea, Philippines
Ozopactus Simon, 1889 – Venezuela
Pachistopelma Pocock, 1901 – Brazil
Pamphobeteus Pocock, 1901 – South America, Panama
Parvicarina Galleti-Lima, Hamilton, Borges & Guadanucci, 2023
Pelinobius Karsch, 1885 – Kenya, Tanzania
Phlogiellus Pocock, 1897 – Asia, Papua New Guinea
Phoneyusa Karsch, 1884 – Africa
Phormictopus Pocock, 1901 – Cuba, Argentina, Brazil
Phormingochilus Pocock, 1895 – Indonesia
Phrixotrichus Simon, 1889 – Chile, Argentina
Plesiopelma Pocock, 1901 – South America
Plesiophrictus Pocock, 1899 – India, Micronesia, Sri Lanka
Poecilotheria Simon, 1885 – Sri Lanka, India
Proshapalopus Mello-Leitão, 1923 – Brazil, Colombia
Psalistops Simon, 1889 - Colombia and Venezuela
Psalmopoeus Pocock, 1895 – Trinidad, South America, Central America, Mexico
Psednocnemis West, Nunn & Hogg, 2012 – Malaysia, Indonesia
Pseudhapalopus Strand, 1907 – South America, Trinidad
Pseudoschizopelma Smith, 1995 - Mexico
Pterinochilus Pocock, 1897 – Africa
Pterinopelma Pocock, 1901 – Brazil
Reichlingia Rudloff, 2001 – Belize
Reversopelma Schmidt, 2001 – Ecuador or Peru
Sahydroaraneus Mirza & Sanap, 2014 – India
Sandinista Longhorn & Gabriel, 2019
Schismatothele Karsch, 1879 – Venezuela, Colombia
Schizopelma F. O. Pickard-Cambridge, 1897 – Mexico
Scopelobates Simon, 1903 – Dominican Republic
Selenocosmia Ausserer, 1871 – Oceania, Asia
Selenogyrus Pocock, 1897 – Côte d'Ivoire, Sierra Leone
Selenotholus Hogg, 1902 – Australia
Selenotypus Pocock, 1895 – Australia
Sericopelma Ausserer, 1875 – Central America, Brazil, Mexico
Sickius Soares & Camargo, 1948 – Brazil
Sphaerobothria Karsch, 1879 – Costa Rica, Panama
Spinosatibiapalpus Gabriel & Sherwood, 2020
Stichoplastoris Rudloff, 1997 – El Salvador, Costa Rica, Panama
Stromatopelma Karsch, 1881 – Africa
Taksinus Songsangchote, Sippawat, Khaikaew & Chomphuphuang, 2022 - Thailand
Tapinauchenius Ausserer, 1871 – South America, Saint Vincent and the Grenadines
Tekoapora Galleti-Lima, Hamilton, Borges & Guadanucci, 2023
Thalerommata Ausserer, 1875 — Colombia, Mexico
Theraphosa Thorell, 1870 – South America
Thrigmopoeus Pocock, 1899 – India
Thrixopelma Schmidt, 1994 – Peru, Chile
Tliltocatl - Mexico, Costa Rica, Guatemala
Tmesiphantes Simon, 1892 – Brazil
Trichognathella Gallon, 2004 – South Africa
Trichopelma Simon, 1888 – Caribbean, South America, Central America
Typhochlaena C. L. Koch, 1850 – Brazil
Umbyquyra Gargiulo, Brescovit & Lucas, 2018 – Bolivia, Brazil
Urupelma Kaderka, Lüddecke, Řezáč, Řezáčová & Hüsser, 2023
Vitalius Lucas, Silva & Bertani, 1993 – Brazil, Argentina
Xenesthis Simon, 1891 – Panama, Venezuela, Colombia
Yanomamius Bertani & Almeida, 2021 – Brazil, Venezuela
Ybyrapora Fukushima & Bertani, 2017 – Brazil
Former genera:
Ami Pérez-Miles, 2008 → Neischnocolus
Barropelma Chamberlin, 1940 → Neischnocolus
Eurypelmella Strand, 1907, nomen dubium
Magulla Simon, 1892 → Tmesiphantes
Melloleitaoina Gerschman & Schiapelli, 1960 → Tmesiphantes
Fossil record
Fossils of mygalomorph spiders date back to the Triassic. One species assigned to the Theraphosidae is Protertheraphosa spinipes, found in Burmese amber, which is dated to the mid and late Cretaceous.
| Biology and health sciences | Arachnids | null |
18581278 | https://en.wikipedia.org/wiki/Digit%20%28anatomy%29 | Digit (anatomy) | A digit is one of several most distal parts of a limb, such as fingers or toes, present in many vertebrates.
Names
Some languages have different names for hand and foot digits (English: respectively "finger" and "toe", German: "Finger" and "Zeh", French: "doigt" and "orteil").
In other languages, e.g. Arabic, Russian, Polish, Spanish, Portuguese, Italian, Czech, Tagalog, Turkish, Bulgarian, and Persian, there are no specific one-word names for fingers and toes; these are called "digit of the hand" or "digit of the foot" instead. In Japanese, yubi (指) can mean either, depending on context.
Human digits
Humans normally have five digits on each extremity. Each digit is formed by several bones called phalanges, surrounded by soft tissue. Human fingers normally have a nail at the distal phalanx. The phenomenon of polydactyly occurs when extra digits are present; fewer digits than normal are also possible, for instance in ectrodactyly. Whether such a mutation can be surgically corrected, and whether such correction is indicated, is case-dependent. For instance the former chess world champion Mikhail Tal lived all his life with only three right-hand fingers.
Brain representation
Each finger has an orderly somatotopic representation on the cerebral cortex in the somatosensory cortex area 3b, part of area 1 and a distributed, overlapping representation in the supplementary motor area and primary motor area.
The somatosensory cortex representation of the hand is a dynamic reflection of the fingers on the external hand: in syndactyly people have a clubhand of webbed, shortened fingers. However, not only are the fingers of their hands fused, but the cortical maps of their individual fingers also form a club hand. The fingers can be surgically divided to make a more useful hand. Surgeons did this at the Institute of Reconstructive Plastic Surgery in New York to a 32-year-old man with the initials O. G.. They touched O. G.’s fingers before and after surgery while using MRI brain scans. Before the surgery, the fingers mapped onto his brain were fused close together; afterward, the maps of his individual fingers did indeed separate and take the layout corresponding to a normal hand.
Evolution
Two ideas about the homology of arms, hands, and digits exist.
That digits are unique to tetrapods; and
That homologous antecedents to digits were present in the fins of early sarcopterygian fish.
Until recently, few transitional forms were known to elaborate on this transition. One particular example is Panderichthys, a coastal fish from the Devonian period 385 million years ago. Prior to 2008, Panderichthys was interpreted as having a fin terminating at a single large plate surrounded by lepidotrichia (fin rays). However, a 2008 study by Boisvert et al. determined that this was mistaken. They discovered that the final bony portion of the fin in Panderichthys is split into at least four fin radials, bones similar to rudimentary fingers.
Thus, in the evolution of tetrapods a shift occurred where the outermost rays of the fins were lost and replaced by the inner radials, which evolve into the earliest digits. This change is consistent with additional evidence from the embryology of actinopterygians, sharks and lungfish. Pre-existing distal radials in these modern fish develop in a very similar way to the digits of tetrapods.
Several rows of digit-like distal fin radials are present in Tiktaalik, a much more complete Devonian vertebrate described in 2006. Though frequently described as the missing link between fishes and tetrapods, the exact relationship between Panderichthys, Tiktaalik, and tetrapods are yet to be fully resolved. Tiktaalik had some features of the forefin more similar to earlier fish, such as a large ulnare and a distinct axis of larger bones down the middle of the fin. According to Boisvert et al. (2008), "It is difficult to say whether this character distribution implies that Tiktaalik is autapomorphic, that Panderichthys and tetrapods are convergent, or that Panderichthys is closer to tetrapods than Tiktaalik. At any rate, it demonstrates that the fish–tetrapod transition was accompanied by significant character incongruence in functionally important structures."p. 638.
Digit-like radials are also known in the rhizodont fish Sauripterus, though this is likely a case of convergent evolution. Elpistostege, a tetrapodomorph fish closely related to Tiktaalik, preserves one of the most tetrapod-like hands in any prehistoric fish. The hand of Elpisostege had 19 distal fin radials arranged into blocks up to four radials long. These sequential blocks of radials are very similar to digits.
Bird and theropod dinosaur digits
Birds and theropod dinosaurs (from which birds evolved) have three digits on their hands. Paradoxically the two digits that are missing are different: the bird hand (embedded in the wing) is thought to derive from the second, third and fourth digits of the ancestral five-digit hand. In contrast, the theropod dinosaurs seem to have the first, second and third digits. Recently a Jurassic theropod intermediate fossil Limusaurus has been found in the Junggar Basin in western China that has a complex mix: it has a first digit stub and full second, third and fourth digits but its wrist bones are like those that are associated with the second, third and fourth digits while its finger bones are those of the first, second and third digits. This suggests the evolution of digits in birds resulted from a "shift in digit identity [that] characterized early stages of theropod evolution"
| Biology and health sciences | External anatomy and regions of the body | Biology |
18581424 | https://en.wikipedia.org/wiki/Silverfish | Silverfish | The silverfish (Lepisma saccharinum) is a species of small, primitive, wingless insect in the order Zygentoma (formerly Thysanura). Its common name derives from the insect's silvery light grey colour, combined with the fish-like appearance of its movements. The scientific name (L. saccharinum) indicates that the silverfish's diet consists of carbohydrates such as sugar or starches. While the common name silverfish is used throughout the global literature to refer to various species of Zygentoma, the Entomological Society of America restricts use of the term solely for Lepisma saccharinum.
Description
The silverfish is a nocturnal insect typically long. Its abdomen tapers at the end, giving it a fish-like appearance. The newly hatched are whitish, but develop a greyish hue and metallic sheen as they get older. It has two long cerci and one terminal filament at the tip of the abdomen between the cerci. It also has two small compound eyes, although other members of Zygentoma are eyeless, such as the family Nicoletiidae.
The silverfish, like other species in Apterygota, is wingless. It has long antennae, and moves in a wiggling motion that resembles the movement of a fish. This, coupled with its appearance and silvery scales, inspires its common name. Silverfish can regenerate lost terminal filaments and antennae within four weeks. Silverfish typically live for up to three years.
The silverfish is an agile runner. It avoids light.
Distribution
Silverfish are a cosmopolitan species, found in Africa, the Americas, Australia, Eurasia, and parts of the Pacific. They inhabit moist areas, requiring a relative humidity between 75% and 95%. In urban areas, they can be found in attics, basements, bathtubs, sinks, kitchens, old books, classrooms, and showers.
Reproduction and life cycle
Before silverfish reproduce, they carry out a ritual involving three phases, which may last over half an hour. In the first phase, the male and female stand face to face, their vibrating antennae touching, then repeatedly back off and return to this position. In the second phase, the male runs away and the female chases him. In the third phase, the male and female stand side by side and head to tail, with the male vibrating his tail against the female. Finally, the male lays a spermatophore, a sperm capsule covered in gossamer, which the female takes into her body via her ovipositor to fertilize her eggs. The female lays groups of fewer than 60 eggs at once, deposited in small crevices. The eggs are oval-shaped, whitish, about long, and take between two weeks and two months to hatch. A silverfish usually lays fewer than 100 eggs in her lifetime.
When the nymphs hatch, they are whitish in colour, and look like smaller adults. As they moult, young silverfish develop a greyish appearance and a metallic sheen, eventually becoming adults after three months to three years. They may go through 17 to 66 moults in their lifetimes, sometimes 30 in a single year—many more than most insects. Silverfish are among the few types of insect that continue to moult after reaching adulthood,
with an estimated lifespan of around 2 to 8 years.
Ecology
Silverfish are able to digest cellulose by themselves, thanks to the cellulase produced by their midgut. They consume matter that contains polysaccharides, such as starches and dextrin in adhesives. These include book bindings, carpet, clothing, coffee, dandruff, glue, hair, some paints, paper, photos, plaster, and sugar. They will damage wallpaper in order to consume the paste. Silverfish can also cause damage to tapestries. Other substances they may eat include cotton, dead insects, linen, silk, leftover crumbs, or even their own exuviae (moulted exoskeleton). During famine, a silverfish may even consume leather and synthetic fabrics. Silverfish can live for a year or more without eating if water is available.
Silverfish are considered household pests, due to their consumption and destruction of property. However, although they are responsible for the contamination of food and other types of damage, they do not transmit disease. House centipedes and spiders such as the spitting spider Scytodes thoracica are known to be predators of silverfish.
The essential oil of the Japanese cedar Cryptomeria japonica has been investigated as a repellent and insecticide against L. saccharinum, with promising results: filter paper impregnated with oil repelled 80% of silverfish at a gas concentration of 0.01 mg/cm, and an exposure of 0.16 mg/cm for 10 hours caused a 100% mortality rate.
Etymology and nomenclature
The scientific name for the species is Lepisma saccharinum (originally saccharina; Linnaeus' 1758 description here), due to its tendency to eat starchy foods high in carbohydrates and protein, such as dextrin. However, the insect's more common name comes from its distinctive metallic appearance and fish-like shape. While the scientific name was established by Carl Linnaeus in his 1758 10th edition of Systema Naturae, the common name has been in use since at least 1855. Most authors have historically treated the nomenclatural gender of Lepisma as feminine (also as specified in ICZN Direction 71 issued in 1957), but in 2018 the International Commission on Zoological Nomenclature reversed this decision, issuing a new formal ruling (ICZN Opinion 2427) stating the gender of Lepisma (and all genera with that ending) is neuter, following ICZN Article 30, which resulted in changes to the spelling of several well-known species, including Lepisma saccharinum.
Evolution
The predecessors of silverfish, along with those of jumping bristletails, are considered the earliest and most primitive insects. They evolved at the latest in mid-Devonian and possibly as early as late Silurian more than 400 million years ago. Some fossilized arthropod trackways from the Paleozoic Era, known as Stiaria intermedia and often attributed to jumping bristletails, may have been produced by silverfish.
Similar species
Other similar insect species are also known as silverfish. Two other silverfish are common in North America, Ctenolepisma longicaudatum and Ctenolepisma quadriseriatum. Ctenolepisma urbanum is known as the urban silverfish.
The Australian species most commonly referred to as silverfish is a different lepismatid, Acrotelsella devriesiana. The firebrat (Thermobia domestica) is like a silverfish, but with a mottled gray and brown body.
| Biology and health sciences | Zygentoma | Animals |
18582186 | https://en.wikipedia.org/wiki/Benzene | Benzene | Benzene is an organic chemical compound with the molecular formula C6H6. The benzene molecule is composed of six carbon atoms joined in a planar hexagonal ring with one hydrogen atom attached to each. Because it contains only carbon and hydrogen atoms, benzene is classed as a hydrocarbon.
Benzene is a natural constituent of petroleum and is one of the elementary petrochemicals. Due to the cyclic continuous pi bonds between the carbon atoms, benzene is classed as an aromatic hydrocarbon. Benzene is a colorless and highly flammable liquid with a sweet smell, and is partially responsible for the aroma of gasoline. It is used primarily as a precursor to the manufacture of chemicals with more complex structures, such as ethylbenzene and cumene, of which billions of kilograms are produced annually. Although benzene is a major industrial chemical, it finds limited use in consumer items because of its toxicity. Benzene is a volatile organic compound.
Benzene is classified as a carcinogen. Its particular effects on human health, such as the long-term results of accidental exposure, have been reported on by news organizations such as The New York Times. For instance, a 2022 article stated that benzene contamination in the Boston metropolitan area caused hazardous conditions in multiple places, with the publication noting that the compound may eventually cause leukemia in some individuals.
History
Discovery
The word "benzene" derives from "gum benzoin" (benzoin resin), an aromatic resin known since ancient times in Southeast Asia, and later to European pharmacists and perfumers in the 16th century via trade routes. An acidic material was derived from benzoin by sublimation, and named "flowers of benzoin", or benzoic acid. The hydrocarbon derived from benzoic acid thus acquired the name benzin, benzol, or benzene. Michael Faraday first isolated and identified benzene in 1825 from the oily residue derived from the production of illuminating gas, giving it the name bicarburet of hydrogen. In 1833, Eilhard Mitscherlich produced it by distilling benzoic acid (from gum benzoin) and lime. He gave the compound the name benzin. In 1836, the French chemist Auguste Laurent named the substance "phène"; this word has become the root of the English word "phenol", which is hydroxylated benzene, and "phenyl", the radical formed by abstraction of a hydrogen atom from benzene.
In 1845, Charles Blachford Mansfield, working under August Wilhelm von Hofmann, isolated benzene from coal tar. Four years later, Mansfield began the first industrial-scale production of benzene, based on the coal-tar method. Gradually, the sense developed among chemists that a number of substances were chemically related to benzene, comprising a diverse chemical family. In 1855, Hofmann was the first to apply the word "aromatic" to designate this family relationship, after a characteristic property of many of its members. In 1997, benzene was detected in deep space.
Ring formula
The empirical formula for benzene was long known, but its highly polyunsaturated structure, with just one hydrogen atom for each carbon atom, was challenging to determine. Archibald Scott Couper in 1858 and Johann Josef Loschmidt in 1861 suggested possible structures that contained multiple double bonds or multiple rings, but in these years very little was known about aromatic chemistry, and so chemists were unable to adduce appropriate evidence to favor any particular formula.
But many chemists had begun to work on aromatic substances, especially in Germany, and relevant data was coming fast. In 1865, the German chemist Friedrich August Kekulé published a paper in French (for he was then teaching in Francophone Belgium) suggesting that the structure contained a ring of six carbon atoms with alternating single and double bonds. The next year he published a much longer paper in German on the same subject. Kekulé used evidence that had accumulated in the intervening years—namely, that there always appeared to be only one isomer of any monoderivative of benzene, and that there always appeared to be exactly three isomers of every disubstituted derivative—now understood to correspond to the ortho, meta, and para patterns of arene substitution—to argue in support of his proposed structure. Kekulé's symmetrical ring could explain these curious facts, as well as benzene's 1:1 carbon-hydrogen ratio.
The new understanding of benzene, and hence of all aromatic compounds, proved to be so important for both pure and applied chemistry that in 1890 the German Chemical Society organized an elaborate appreciation in Kekulé's honor, celebrating the twenty-fifth anniversary of his first benzene paper. Here Kekulé spoke of the creation of the theory. He said that he had discovered the ring shape of the benzene molecule after having a reverie or day-dream of a snake biting its own tail (a symbol in ancient cultures known as the ouroboros). This vision, he said, came to him after years of studying the nature of carbon-carbon bonds. This was seven years after he had solved the problem of how carbon atoms could bond to up to four other atoms at the same time. Curiously, a similar, humorous depiction of benzene had appeared in 1886 in a pamphlet entitled Berichte der Durstigen Chemischen Gesellschaft (Journal of the Thirsty Chemical Society), a parody of the Berichte der Deutschen Chemischen Gesellschaft, only the parody had monkeys seizing each other in a circle, rather than snakes as in Kekulé's anecdote. Some historians have suggested that the parody was a lampoon of the snake anecdote, possibly already well known through oral transmission even if it had not yet appeared in print. Kekulé's 1890 speech in which this anecdote appeared has been translated into English. If the anecdote is the memory of a real event, circumstances mentioned in the story suggest that it must have happened early in 1862.
In 1929, the cyclic nature of benzene was finally confirmed by the crystallographer Kathleen Lonsdale using X-ray diffraction methods. Using large crystals of hexamethylbenzene, a benzene derivative with the same core of six carbon atoms, Lonsdale obtained diffraction patterns. Through calculating more than thirty parameters, Lonsdale demonstrated that the benzene ring could not be anything but a flat hexagon, and provided accurate distances for all carbon-carbon bonds in the molecule.
Nomenclature
The German chemist Wilhelm Körner suggested the prefixes ortho-, meta-, para- to distinguish di-substituted benzene derivatives in 1867; however, he did not use the prefixes to distinguish the relative positions of the substituents on a benzene ring. It was the German chemist Carl Gräbe who, in 1869, first used the prefixes ortho-, meta-, para- to denote specific relative locations of the substituents on a di-substituted aromatic ring (viz, naphthalene). In 1870, the German chemist Viktor Meyer first applied Gräbe's nomenclature to benzene.
Early applications
In 1903, Ludwig Roselius popularized the use of benzene to decaffeinate coffee. This discovery led to the production of Sanka. This process was later discontinued. Benzene was historically used as a significant component in many consumer products such as liquid wrench, several paint strippers, rubber cements, spot removers, and other products. Manufacture of some of these benzene-containing formulations ceased in about 1950, although Liquid Wrench continued to contain significant amounts of benzene until the late 1970s.
Occurrence
Trace amounts of benzene are found in petroleum and coal. It is a byproduct of the incomplete combustion of many materials. For commercial use, until World War II, much of benzene was obtained as a by-product of coke production (or "coke-oven light oil") for the steel industry. However, in the 1950s, increased demand for benzene, especially from the growing polymers industry, necessitated the production of benzene from petroleum. Today, most benzene comes from the petrochemical industry, with only a small fraction being produced from coal. Benzene has been detected on Mars.
Structure
X-ray diffraction shows that all six carbon-carbon bonds in benzene are of the same length, at 140 picometres (pm). The C–C bond lengths are greater than a double bond (135 pm) but shorter than a single bond (147 pm). This intermediate distance is caused by electron delocalization: the electrons for C=C bonding are distributed equally between each of the six carbon atoms. Benzene has 6 hydrogen atoms, fewer than the corresponding parent alkane, hexane, which has 14. Benzene and cyclohexane have a similar structure, only the ring of delocalized electrons and the loss of one hydrogen per carbon distinguishes it from cyclohexane. The molecule is planar. The molecular orbital description involves the formation of three delocalized π orbitals spanning all six carbon atoms, while the valence bond description involves a superposition of resonance structures. It is likely that this stability contributes to the peculiar molecular and chemical properties known as aromaticity. To reflect the delocalized nature of the bonding, benzene is often depicted with a circle inside a hexagonal arrangement of carbon atoms.
Derivatives of benzene occur sufficiently often as a component of organic molecules, so much so that the Unicode Consortium has allocated a symbol in the Miscellaneous Technical block with the code U+232C (⌬) to represent it with three double bonds, and U+23E3 (⏣) for a delocalized version.
Benzene derivatives
Many important chemical compounds are derived from benzene by replacing one or more of its hydrogen atoms with another functional group. Examples of simple benzene derivatives are phenol, toluene, and aniline, abbreviated PhOH, PhMe, and PhNH2, respectively. Linking benzene rings gives biphenyl, C6H5–C6H5. Further loss of hydrogen gives "fused" aromatic hydrocarbons, such as naphthalene, anthracene, phenanthrene, and pyrene. The limit of the fusion process is the hydrogen-free allotrope of carbon, graphite.
In heterocycles, carbon atoms in the benzene ring are replaced with other elements. The most important variations contain nitrogen. Replacing one CH with N gives the compound pyridine, C5H5N. Although benzene and pyridine are structurally related, benzene cannot be converted into pyridine. Replacement of a second CH bond with N gives, depending on the location of the second N, pyridazine, pyrimidine, or pyrazine.
Production
Four chemical processes contribute to industrial benzene production: catalytic reforming, toluene hydrodealkylation, toluene disproportionation, and steam cracking etc. According to the ATSDR Toxicological Profile for benzene, between 1978 and 1981, catalytic reformates accounted for approximately 44–50% of the total U.S. benzene production.
Catalytic reforming
In catalytic reforming, a mixture of hydrocarbons with boiling points between 60 and 200 °C is blended with hydrogen gas and then exposed to a bifunctional platinum chloride or rhenium chloride catalyst at 500–525 °C and pressures ranging from 8–50 atm. Under these conditions, aliphatic hydrocarbons form rings and lose hydrogen to become aromatic hydrocarbons. The aromatic products of the reaction are then separated from the reaction mixture (or reformate) by extraction with any one of a number of solvents, including diethylene glycol or sulfolane, and benzene is then separated from the other aromatics by distillation. The extraction step of aromatics from the reformate is designed to produce aromatics with lowest non-aromatic components. Recovery of the aromatics, commonly referred to as BTX (benzene, toluene and xylene isomers), involves such extraction and distillation steps.
In similar fashion to this catalytic reforming, UOP and BP commercialized a method from LPG (mainly propane and butane) to aromatics.
Toluene hydrodealkylation
Toluene hydrodealkylation converts toluene to benzene. In this hydrogen-intensive process, toluene is mixed with hydrogen, then passed over a chromium, molybdenum, or platinum oxide catalyst at 500–650 °C and 20–60 atm pressure. Sometimes, higher temperatures are used instead of a catalyst (at the similar reaction condition). Under these conditions, toluene undergoes dealkylation to benzene and methane:
This irreversible reaction is accompanied by an equilibrium side reaction that produces
biphenyl (diphenyl) at higher temperature:
2 +
If the raw material stream contains much non-aromatic components (paraffins or naphthenes), those are likely decomposed to lower hydrocarbons such as methane, which increases the consumption of hydrogen.
A typical reaction yield exceeds 95%. Sometimes, xylenes and heavier aromatics are used in place of toluene, with similar efficiency.
This is often called "on-purpose" methodology to produce benzene, compared to conventional BTX (benzene-toluene-xylene) extraction processes.
Toluene disproportionation
Toluene disproportionation (TDP) is the conversion of toluene to benzene and xylene.
Given that demand for para-xylene (p-xylene) substantially exceeds demand for other xylene isomers, a refinement of the TDP process called Selective TDP (STDP) may be used. In this process, the xylene stream exiting the TDP unit is approximately 90% p-xylene. In some systems, even the benzene-to-xylenes ratio is modified to favor xylenes.
Steam cracking
Steam cracking is the process for producing ethylene and other alkenes from aliphatic hydrocarbons. Depending on the feedstock used to produce the olefins, steam cracking can produce a benzene-rich liquid by-product called pyrolysis gasoline. Pyrolysis gasoline can be blended with other hydrocarbons as a gasoline additive, or routed through an extraction process to recover BTX aromatics (benzene, toluene and xylenes).
Other methods
Although of no commercial significance, many other routes to benzene exist. Phenol and halobenzenes can be reduced with metals. Benzoic acid and its salts undergo decarboxylation to benzene. The reaction of the diazonium compound derived from aniline with hypophosphorus acid gives benzene. Alkyne trimerisation of acetylene gives benzene. Complete decarboxylation of mellitic acid gives benzene.
Uses
Benzene is used mainly as an intermediate to make other chemicals, above all ethylbenzene (and other alkylbenzenes), cumene, cyclohexane, and nitrobenzene. In 1988 it was reported that two-thirds of all chemicals on the American Chemical Society's lists contained at least one benzene ring. More than half of the entire benzene production is processed into ethylbenzene, a precursor to styrene, which is used to make polymers and plastics like polystyrene. Some 20% of the benzene production is used to manufacture cumene, which is needed to produce phenol and acetone for resins and adhesives. Cyclohexane consumes around 10% of the world's benzene production; it is primarily used in the manufacture of nylon fibers, which are processed into textiles and engineering plastics. Smaller amounts of benzene are used to make some types of rubbers, lubricants, dyes, detergents, drugs, explosives, and pesticides. In 2013, the biggest consumer country of benzene was China, followed by the USA. Benzene production is currently expanding in the Middle East and in Africa, whereas production capacities in Western Europe and North America are stagnating.
Toluene is now often used as a substitute for benzene, for instance as a fuel additive. The solvent-properties of the two are similar, but toluene is less toxic and has a wider liquid range. Toluene is also processed into benzene.
Component of gasoline
As a gasoline (petrol) additive, benzene increases the octane rating and reduces knocking. As a consequence, gasoline often contained several percent benzene before the 1950s, when tetraethyl lead replaced it as the most widely used antiknock additive. With the global phaseout of leaded gasoline, benzene has made a comeback as a gasoline additive in some nations. In the United States, concern over its negative health effects and the possibility of benzene entering the groundwater has led to stringent regulation of gasoline's benzene content, with limits typically around 1%. European petrol specifications now contain the same 1% limit on benzene content. The United States Environmental Protection Agency introduced new regulations in 2011 that lowered the benzene content in gasoline to 0.62%.
In some European languages, the word for petroleum or gasoline is an exact cognate of "benzene". For instance in Catalan the word 'benzina' can be used for gasoline, though now it is relatively rare.
Reactions
The most common reactions of benzene involve substitution of a proton by other groups. Electrophilic aromatic substitution is a general method of derivatizing benzene. Benzene is sufficiently nucleophilic that it undergoes substitution by acylium ions and alkyl carbocations to give substituted derivatives.
The most widely practiced example of this reaction is the ethylation of benzene.
Approximately 24,700,000 tons were produced in 1999. Highly instructive but of far less industrial significance is the Friedel-Crafts alkylation of benzene (and many other aromatic rings) using an alkyl halide in the presence of a strong Lewis acid catalyst. Similarly, the Friedel-Crafts acylation is a related example of electrophilic aromatic substitution. The reaction involves the acylation of benzene (or many other aromatic rings) with an acyl chloride using a strong Lewis acid catalyst such as aluminium chloride or Iron(III) chloride.
Sulfonation, chlorination, nitration
Using electrophilic aromatic substitution, many functional groups are introduced onto the benzene framework. Sulfonation of benzene involves the use of oleum, a mixture of sulfuric acid with sulfur trioxide. Sulfonated benzene derivatives are useful detergents. In nitration, benzene reacts with nitronium ions (NO2+), which is a strong electrophile produced by combining sulfuric and nitric acids. Nitrobenzene is the precursor to aniline. Chlorination is achieved with chlorine to give chlorobenzene in the presence of a Lewis acid catalyst such as aluminium tri-chloride.
Hydrogenation
Via hydrogenation, benzene and its derivatives convert to cyclohexane and derivatives. This reaction is achieved by the use of high pressures of hydrogen in the presence of heterogeneous catalysts, such as finely divided nickel. Whereas alkenes can be hydrogenated near room temperatures, benzene and related compounds are more reluctant substrates, requiring temperatures >100 °C. This reaction is practiced on a large scale industrially. In the absence of the catalyst, benzene is impervious to hydrogen. Hydrogenation cannot be stopped to give cyclohexene or cyclohexadienes as these are superior substrates. Birch reduction, a non catalytic process, however selectively hydrogenates benzene to the diene.
Metal complexes
Benzene is an excellent ligand in the organometallic chemistry of low-valent metals. Important examples include the sandwich and half-sandwich complexes, respectively, Cr(C6H6)2 and [RuCl2(C6H6)]2.
Health effects
Benzene is classified as a carcinogen, which increases the risk of cancer and other illnesses, and is also a notorious cause of bone marrow failure. Substantial quantities of epidemiologic, clinical, and laboratory data link benzene to aplastic anemia, acute leukemia, bone marrow abnormalities and cardiovascular disease. The specific hematologic malignancies that benzene is associated with include: acute myeloid leukemia (AML), aplastic anemia, myelodysplastic syndrome (MDS), acute lymphoblastic leukemia (ALL), and chronic myeloid leukemia (CML).
Carcinogenic activity of benzene was discovered by Swedish pharmacologist in 1897 on female workers of a tire-making factory. The American Petroleum Institute (API) stated in 1948 that "it is generally considered that the only absolutely safe concentration for benzene is zero". There is no safe exposure level; even tiny amounts can cause harm. The US Department of Health and Human Services (DHHS) classifies benzene as a human carcinogen. Long-term exposure to excessive levels of benzene in the air causes leukemia, a potentially fatal cancer of the blood-forming organs. In particular, acute myeloid leukemia or acute nonlymphocytic leukemia (AML & ANLL) is caused by benzene. IARC rated benzene as "known to be carcinogenic to humans" (Group 1).
As benzene is ubiquitous in gasoline and hydrocarbon fuels that are in use everywhere, human exposure to benzene is a global health problem. Benzene targets the liver, kidney, lung, heart and brain and can cause DNA strand breaks and chromosomal damage, hence is teratogenic and mutagenic. Benzene causes cancer in animals including humans. Benzene has been shown to cause cancer in both sexes of multiple species of laboratory animals exposed via various routes.
Exposure to benzene
According to the Agency for Toxic Substances and Disease Registry (ATSDR) (2007), benzene is both a synthetically made and naturally occurring chemical from processes that include: volcanic eruptions, wild fires, synthesis of chemicals such as phenol, production of synthetic fibers, and fabrication of rubbers, lubricants, pesticides, medications, and dyes. The major sources of benzene exposure are tobacco smoke, automobile service stations, exhaust from motor vehicles, and industrial emissions; however, ingestion and dermal absorption of benzene can also occur through contact with contaminated water. Benzene is hepatically metabolized and excreted in the urine. Measurement of air and water levels of benzene is accomplished through collection via activated charcoal tubes, which are then analyzed with a gas chromatograph. The measurement of benzene in humans can be accomplished via urine, blood, and breath tests; however, all of these have their limitations because benzene is rapidly metabolized in the human body.
Exposure to benzene may lead progressively to aplastic anemia, leukaemia, and multiple myeloma.
OSHA regulates levels of benzene in the workplace. The maximum allowable amount of benzene in workroom air during an 8-hour workday, 40-hour workweek is 1 ppm. As benzene can cause cancer, NIOSH recommends that all workers wear special breathing equipment when they are likely to be exposed to benzene at levels exceeding the recommended (8-hour) exposure limit of 0.1 ppm.
Benzene exposure limits
The United States Environmental Protection Agency has set a maximum contaminant level for benzene in drinking water at 0.005 mg/L (5 ppb), as promulgated via the U.S. National Primary Drinking Water Regulations. This regulation is based on preventing benzene leukemogenesis. The maximum contaminant level goal (MCLG), a nonenforceable health goal that would allow an adequate margin of safety for the prevention of adverse effects, is zero benzene concentration in drinking water. The EPA requires that spills or accidental releases into the environment of 10 pounds (4.5 kg) or more of benzene be reported.
The U.S. Occupational Safety and Health Administration (OSHA) has set a permissible exposure limit of 1 part of benzene per million parts of air (1 ppm) in the workplace during an 8-hour workday, 40-hour workweek. The short term exposure limit for airborne benzene is 5 ppm for 15 minutes. These legal limits were based on studies demonstrating compelling evidence of health risk to workers exposed to benzene. The risk from exposure to 1 ppm for a working lifetime has been estimated as 5 excess leukemia deaths per 1,000 employees exposed. (This estimate assumes no threshold for benzene's carcinogenic effects.) OSHA has also established an action level of 0.5 ppm to encourage even lower exposures in the workplace.
The U.S. National Institute for Occupational Safety and Health (NIOSH) revised the Immediately Dangerous to Life and Health (IDLH) concentration for benzene to 500 ppm. The current NIOSH definition for an IDLH condition, as given in the NIOSH Respirator Selection Logic, is one that poses a threat of exposure to airborne contaminants when that exposure is likely to cause death or immediate or delayed permanent adverse health effects or prevent escape from such an environment. The purpose of establishing an IDLH value is (1) to ensure that the worker can escape from a given contaminated environment in the event of failure of the respiratory protection equipment and (2) is considered a maximum level above which only a highly reliable breathing apparatus providing maximum worker protection is permitted. In September 1995, NIOSH issued a new policy for developing recommended exposure limits (RELs) for substances, including carcinogens. As benzene can cause cancer, NIOSH recommends that all workers wear special breathing equipment when they are likely to be exposed to benzene at levels exceeding the REL (10-hour) of 0.1 ppm. The NIOSH short-term exposure limit (STEL – 15 min) is 1 ppm.
American Conference of Governmental Industrial Hygienists (ACGIH) adopted Threshold Limit Values (TLVs) for benzene at 0.5 ppm TWA and 2.5 ppm STEL.
Toxicology
Biomarkers of exposure
Several tests can determine exposure to benzene. Benzene itself can be measured in breath, blood or urine, but such testing is usually limited to the first 24 hours post-exposure due to the relatively rapid removal of the chemical by exhalation or biotransformation. Most people in developed countries have measureable baseline levels of benzene and other aromatic petroleum hydrocarbons in their blood. In the body, benzene is enzymatically converted to a series of oxidation products including muconic acid, phenylmercapturic acid, phenol, catechol, hydroquinone and 1,2,4-trihydroxybenzene. Most of these metabolites have some value as biomarkers of human exposure, since they accumulate in the urine in proportion to the extent and duration of exposure, and they may still be present for some days after exposure has ceased. The current ACGIH biological exposure limits for occupational exposure are 500 μg/g creatinine for muconic acid and 25 μg/g creatinine for phenylmercapturic acid in an end-of-shift urine specimen.
Biotransformations
Even if it is not a common substrate for metabolism, benzene can be oxidized by both bacteria and eukaryotes. In bacteria, dioxygenase enzyme can add an oxygen to the ring, and the unstable product is immediately reduced (by NADH) to a cyclic diol with two double bonds, breaking the aromaticity. Next, the diol is newly reduced by NADH to catechol. The catechol is then metabolized to acetyl CoA and succinyl CoA, used by organisms mainly in the citric acid cycle for energy production.
The pathway for the metabolism of benzene is complex and begins in the liver. Several enzymes are involved. These include cytochrome P450 2E1 (CYP2E1), quinine oxidoreductase (NQ01 or DT-diaphorase or NAD(P)H dehydrogenase (quinone 1)), GSH, and myeloperoxidase (MPO). CYP2E1 is involved at multiple steps: converting benzene to oxepin (benzene oxide), phenol to hydroquinone, and hydroquinone to both benzenetriol and catechol. Hydroquinone, benzenetriol and catechol are converted to polyphenols. In the bone marrow, MPO converts these polyphenols to benzoquinones. These intermediates and metabolites induce genotoxicity by multiple mechanisms including inhibition of topoisomerase II (which maintains chromosome structure), disruption of microtubules (which maintains cellular structure and organization), generation of oxygen free radicals (unstable species) that may lead to point mutations, increasing oxidative stress, inducing DNA strand breaks, and altering DNA methylation (which can affect gene expression). NQ01 and GSH shift metabolism away from toxicity. NQ01 metabolizes benzoquinone toward polyphenols (counteracting the effect of MPO). GSH is involved with the formation of phenylmercapturic acid.
Genetic polymorphisms in these enzymes may induce loss of function or gain of function. For example, mutations in CYP2E1 increase activity and result in increased generation of toxic metabolites. NQ01 mutations result in loss of function and may result in decreased detoxification. Myeloperoxidase mutations result in loss of function and may result in decreased generation of toxic metabolites. GSH mutations or deletions result in loss of function and result in decreased detoxification. These genes may be targets for genetic screening for susceptibility to benzene toxicity.
Molecular toxicology
The paradigm of toxicological assessment of benzene is shifting towards the domain of molecular toxicology as it allows understanding of fundamental biological mechanisms in a better way. Glutathione seems to play an important role by protecting against benzene-induced DNA breaks and it is being identified as a new biomarker for exposure and effect. Benzene causes chromosomal aberrations in the peripheral blood leukocytes and bone marrow explaining the higher incidence of leukemia and multiple myeloma caused by chronic exposure. These aberrations can be monitored using fluorescent in situ hybridization (FISH) with DNA probes to assess the effects of benzene along with the hematological tests as markers of hematotoxicity. Benzene metabolism involves enzymes coded for by polymorphic genes. Studies have shown that genotype at these loci may influence susceptibility to the toxic effects of benzene exposure. Individuals carrying variant of NAD(P)H:quinone oxidoreductase 1 (NQO1), microsomal epoxide hydrolase (EPHX) and deletion of the glutathione S-transferase T1 (GSTT1) showed a greater frequency of DNA single-stranded breaks.
Biological oxidation and carcinogenic activity
One way of understanding the carcinogenic effects of benzene is by examining the products of biological oxidation. Pure benzene, for example, oxidizes in the body to produce an epoxide, benzene oxide, which is not excreted readily and can interact with DNA to produce harmful mutations.
Routes of exposure
Inhalation
Outdoor air may contain low levels of benzene from automobile service stations, wood smoke, tobacco smoke, the transfer of gasoline, exhaust from motor vehicles, and industrial emissions. About 50% of the entire nationwide (United States) exposure to benzene results from smoking tobacco or from exposure to tobacco smoke. After smoking 32 cigarettes per day, the smoker would take in about 1.8 milligrams (mg) of benzene. This amount is about 10 times the average daily intake of benzene by nonsmokers.
Inhaled benzene is primarily expelled unchanged through exhalation. In a human study 16.4 to 41.6% of retained benzene was eliminated through the lungs within five to seven hours after a two- to three-hour exposure to 47 to 110 ppm and only 0.07 to 0.2% of the remaining benzene was excreted unchanged in the urine. After exposure to 63 to 405 mg/m3 of benzene for 1 to 5 hours, 51 to 87% was excreted in the urine as phenol over a period of 23 to 50 hours. In another human study, 30% of absorbed dermally applied benzene, which is primarily metabolized in the liver, was excreted as phenol in the urine.
Exposure from soft drinks
Under specific conditions and in the presence of other chemicals benzoic acid (a preservative) and ascorbic acid (Vitamin C) may interact to produce benzene. In March 2006, the official Food Standards Agency in United Kingdom conducted a survey of 150 brands of soft drinks. It found that four contained benzene levels above World Health Organization limits. The affected batches were removed from sale. Similar problems were reported by the FDA in the United States.
Contamination of water supply
In 2005, the water supply to the city of Harbin in China with a population of almost nine million people, was cut off because of a major benzene exposure. Benzene leaked into the Songhua River, which supplies drinking water to the city, after an explosion at a China National Petroleum Corporation (CNPC) factory in the city of Jilin on 13 November 2005.
When plastic water pipes are subject to high heat, the water may be contaminated with benzene.
Genocide
The Nazi German government used benzene administered via injection as one of their many methods for killing.
| Physical sciences | Hydrocarbons | null |
18582230 | https://en.wikipedia.org/wiki/Methane | Methane | Methane ( , ) is a chemical compound with the chemical formula (one carbon atom bonded to four hydrogen atoms). It is a group-14 hydride, the simplest alkane, and the main constituent of natural gas. The abundance of methane on Earth makes it an economically attractive fuel, although capturing and storing it is difficult because it is a gas at standard temperature and pressure. In the Earth's atmosphere methane is transparent to visible light but absorbs infrared radiation, acting as a greenhouse gas. Methane is an organic compound, and among the simplest of organic compounds. Methane is also a hydrocarbon.
Naturally occurring methane is found both below ground and under the seafloor and is formed by both geological and biological processes. The largest reservoir of methane is under the seafloor in the form of methane clathrates. When methane reaches the surface and the atmosphere, it is known as atmospheric methane.
The Earth's atmospheric methane concentration has increased by about 160% since 1750, with the overwhelming percentage caused by human activity. It accounted for 20% of the total radiative forcing from all of the long-lived and globally mixed greenhouse gases, according to the 2021 Intergovernmental Panel on Climate Change report. Strong, rapid and sustained reductions in methane emissions could limit near-term warming and improve air quality by reducing global surface ozone.
Methane has also been detected on other planets, including Mars, which has implications for astrobiology research.
Properties and bonding
Methane is a tetrahedral molecule with four equivalent C–H bonds. Its electronic structure is described by four bonding molecular orbitals (MOs) resulting from the overlap of the valence orbitals on C and H. The lowest-energy MO is the result of the overlap of the 2s orbital on carbon with the in-phase combination of the 1s orbitals on the four hydrogen atoms. Above this energy level is a triply degenerate set of MOs that involve overlap of the 2p orbitals on carbon with various linear combinations of the 1s orbitals on hydrogen. The resulting "three-over-one" bonding scheme is consistent with photoelectron spectroscopic measurements.
Methane is an odorless, colourless and transparent gas at standard temperature and pressure. It does absorb visible light, especially at the red end of the spectrum, due to overtone bands, but the effect is only noticeable if the light path is very long. This is what gives Uranus and Neptune their blue or bluish-green colors, as light passes through their atmospheres containing methane and is then scattered back out.
The familiar smell of natural gas as used in homes is achieved by the addition of an odorant, usually blends containing tert-butylthiol, as a safety measure. Methane has a boiling point of −161.5 °C at a pressure of one atmosphere. As a gas, it is flammable over a range of concentrations (5.4%–17%) in air at standard pressure.
Solid methane exists in several modifications, of which nine are known. Cooling methane at normal pressure results in the formation of methane I. This substance crystallizes in the cubic system (space group Fmm). The positions of the hydrogen atoms are not fixed in methane I, i.e. methane molecules may rotate freely. Therefore, it is a plastic crystal.
Chemical reactions
The primary chemical reactions of methane are combustion, steam reforming to syngas, and halogenation. In general, methane reactions are difficult to control.
Selective oxidation
Partial oxidation of methane to methanol (CH3OH), a more convenient, liquid fuel, is challenging because the reaction typically progresses all the way to carbon dioxide and water even with an insufficient supply of oxygen. The enzyme methane monooxygenase produces methanol from methane, but cannot be used for industrial-scale reactions. Some homogeneously catalyzed systems and heterogeneous systems have been developed, but all have significant drawbacks. These generally operate by generating protected products which are shielded from overoxidation. Examples include the Catalytica system, copper zeolites, and iron zeolites stabilizing the alpha-oxygen active site.
One group of bacteria catalyze methane oxidation with nitrite as the oxidant in the absence of oxygen, giving rise to the so-called anaerobic oxidation of methane.
Acid–base reactions
Like other hydrocarbons, methane is an extremely weak acid. Its pKa in DMSO is estimated to be 56. It cannot be deprotonated in solution, but the conjugate base is known in forms such as methyllithium.
A variety of positive ions derived from methane have been observed, mostly as unstable species in low-pressure gas mixtures. These include methenium or methyl cation , methane cation , and methanium or protonated methane . Some of these have been detected in outer space. Methanium can also be produced as diluted solutions from methane with superacids. Cations with higher charge, such as and , have been studied theoretically and conjectured to be stable.
Despite the strength of its C–H bonds, there is intense interest in catalysts that facilitate C–H bond activation in methane (and other lower numbered alkanes).
Combustion
Methane's heat of combustion is 55.5 MJ/kg. Combustion of methane is a multiple step reaction summarized as follows:
(ΔH = −891 kJ/mol, at standard conditions)
Peters four-step chemistry is a systematically reduced four-step chemistry that explains the burning of methane.
Methane radical reactions
Given appropriate conditions, methane reacts with halogen radicals as follows:
where X is a halogen: fluorine (F), chlorine (Cl), bromine (Br), or iodine (I). This mechanism for this process is called free radical halogenation. It is initiated when UV light or some other radical initiator (like peroxides) produces a halogen atom. A two-step chain reaction ensues in which the halogen atom abstracts a hydrogen atom from a methane molecule, resulting in the formation of a hydrogen halide molecule and a methyl radical (). The methyl radical then reacts with a molecule of the halogen to form a molecule of the halomethane, with a new halogen atom as byproduct. Similar reactions can occur on the halogenated product, leading to replacement of additional hydrogen atoms by halogen atoms with dihalomethane, trihalomethane, and ultimately, tetrahalomethane structures, depending upon reaction conditions and the halogen-to-methane ratio.
This reaction is commonly used with chlorine to produce dichloromethane and chloroform via chloromethane. Carbon tetrachloride can be made with excess chlorine.
Uses
Methane may be transported as a refrigerated liquid (liquefied natural gas, or LNG). While leaks from a refrigerated liquid container are initially heavier than air due to the increased density of the cold gas, the gas at ambient temperature is lighter than air. Gas pipelines distribute large amounts of natural gas, of which methane is the principal component.
Fuel
Methane is used as a fuel for ovens, homes, water heaters, kilns, automobiles, turbines, etc.
As the major constituent of natural gas, methane is important for electricity generation by burning it as a fuel in a gas turbine or steam generator. Compared to other hydrocarbon fuels, methane produces less carbon dioxide for each unit of heat released. At about 891 kJ/mol, methane's heat of combustion is lower than that of any other hydrocarbon, but the ratio of the heat of combustion (891 kJ/mol) to the molecular mass (16.0 g/mol, of which 12.0 g/mol is carbon) shows that methane, being the simplest hydrocarbon, produces more heat per mass unit (55.7 kJ/g) than other complex hydrocarbons. In many areas with a dense enough population, methane is piped into homes and businesses for heating, cooking, and industrial uses. In this context it is usually known as natural gas, which is considered to have an energy content of 39 megajoules per cubic meter, or 1,000 BTU per standard cubic foot. Liquefied natural gas (LNG) is predominantly methane () converted into liquid form for ease of storage or transport.
Rocket propellant
Refined liquid methane as well as LNG is used as a rocket fuel, when combined with liquid oxygen, as in the TQ-12, BE-4, Raptor, and YF-215 engines. Due to the similarities between methane and LNG such engines are commonly grouped together under the term methalox.
As a liquid rocket propellant, a methane/liquid oxygen combination offers the advantage over kerosene/liquid oxygen combination, or kerolox, of producing small exhaust molecules, reducing coking or deposition of soot on engine components. Methane is easier to store than hydrogen due to its higher boiling point and density, as well as its lack of hydrogen embrittlement. The lower molecular weight of the exhaust also increases the fraction of the heat energy which is in the form of kinetic energy available for propulsion, increasing the specific impulse of the rocket. Compared to liquid hydrogen, the specific energy of methane is lower but this disadvantage is offset by methane's greater density and temperature range, allowing for smaller and lighter tankage for a given fuel mass. Liquid methane has a temperature range (91–112 K) nearly compatible with liquid oxygen (54–90 K). The fuel currently sees use in operational launch vehicles such as Zhuque-2, Vulcan and New Glenn as well as in-development launchers such as Starship, Neutron, and Terran R.
Chemical feedstock
Natural gas, which is mostly composed of methane, is used to produce hydrogen gas on an industrial scale. Steam methane reforming (SMR), or simply known as steam reforming, is the standard industrial method of producing commercial bulk hydrogen gas. More than 50 million metric tons are produced annually worldwide (2013), principally from the SMR of natural gas. Much of this hydrogen is used in petroleum refineries, in the production of chemicals and in food processing. Very large quantities of hydrogen are used in the industrial synthesis of ammonia.
At high temperatures (700–1100 °C) and in the presence of a metal-based catalyst (nickel), steam reacts with methane to yield a mixture of CO and , known as "water gas" or "syngas":
This reaction is strongly endothermic (consumes heat, 206 kJ/mol).
Additional hydrogen is obtained by the reaction of CO with water via the water-gas shift reaction:
This reaction is mildly exothermic (produces heat, −41 kJ/mol).
Methane is also subjected to free-radical chlorination in the production of chloromethanes, although methanol is a more typical precursor.
Hydrogen can also be produced via the direct decomposition of methane, also known as methane pyrolysis, which, unlike steam reforming, produces no greenhouse gases (GHG). The heat needed for the reaction can also be GHG emission free, e.g. from concentrated sunlight, renewable electricity, or burning some of the produced hydrogen. If the methane is from biogas then the process can be a carbon sink. Temperatures in excess of 1200 °C are required to break the bonds of methane to produce hydrogen gas and solid carbon.
However, through the use of a suitable catalyst the reaction temperature can be reduced to between 550 and 900 °C depending on the chosen catalyst. Dozens of catalysts have been tested, including unsupported and supported metal catalysts, carbonaceous and metal-carbon catalysts.
The reaction is moderately endothermic as shown in the reaction equation below.
( 74.8 kJ/mol)
Refrigerant
As a refrigerant, methane has the ASHRAE designation R-50.
Generation
Methane can be generated through geological, biological or industrial routes.
Geological routes
The two main routes for geological methane generation are (i) organic (thermally generated, or thermogenic) and (ii) inorganic (abiotic). Thermogenic methane occurs due to the breakup of organic matter at elevated temperatures and pressures in deep sedimentary strata. Most methane in sedimentary basins is thermogenic; therefore, thermogenic methane is the most important source of natural gas. Thermogenic methane components are typically considered to be relic (from an earlier time). Generally, formation of thermogenic methane (at depth) can occur through organic matter breakup, or organic synthesis. Both ways can involve microorganisms (methanogenesis), but may also occur inorganically. The processes involved can also consume methane, with and without microorganisms.
The more important source of methane at depth (crystalline bedrock) is abiotic. Abiotic means that methane is created from inorganic compounds, without biological activity, either through magmatic processes or via water-rock reactions that occur at low temperatures and pressures, like serpentinization.
Biological routes
Most of Earth's methane is biogenic and is produced by methanogenesis, a form of anaerobic respiration only known to be conducted by some members of the domain Archaea. Methanogens occur in landfills and soils, ruminants (for example, cattle), the guts of termites, and the anoxic sediments below the seafloor and the bottom of lakes.
This multistep process is used by these microorganisms for energy. The net reaction of methanogenesis is:
The final step in the process is catalyzed by the enzyme methyl coenzyme M reductase (MCR).
Wetlands
Wetlands are the largest natural sources of methane to the atmosphere, accounting for approximately 20 – 30% of atmospheric methane. Climate change is increasing the amount of methane released from wetlands due to increased temperatures and altered rainfall patterns. This phenomenon is called wetland methane feedback.
Rice cultivation generates as much as 12% of total global methane emissions due to the long-term flooding of rice fields.
Ruminants
Ruminants, such as cattle, belch methane, accounting for about 22% of the U.S. annual methane emissions to the atmosphere. One study reported that the livestock sector in general (primarily cattle, chickens, and pigs) produces 37% of all human-induced methane. A 2013 study estimated that livestock accounted for 44% of human-induced methane and about 15% of human-induced greenhouse gas emissions. Many efforts are underway to reduce livestock methane production, such as medical treatments and dietary adjustments, and to trap the gas to use its combustion energy.
Seafloor sediments
Most of the subseafloor is anoxic because oxygen is removed by aerobic microorganisms within the first few centimeters of the sediment. Below the oxygen-replete seafloor, methanogens produce methane that is either used by other organisms or becomes trapped in gas hydrates. These other organisms that utilize methane for energy are known as methanotrophs ('methane-eating'), and are the main reason why little methane generated at depth reaches the sea surface. Consortia of Archaea and Bacteria have been found to oxidize methane via anaerobic oxidation of methane (AOM); the organisms responsible for this are anaerobic methanotrophic Archaea (ANME) and sulfate-reducing bacteria (SRB).
Industrial routes
Given its cheap abundance in natural gas, there is little incentive to produce methane industrially. Methane can be produced by hydrogenating carbon dioxide through the Sabatier process. Methane is also a side product of the hydrogenation of carbon monoxide in the Fischer–Tropsch process, which is practiced on a large scale to produce longer-chain molecules than methane.
An example of large-scale coal-to-methane gasification is the Great Plains Synfuels plant, started in 1984 in Beulah, North Dakota as a way to develop abundant local resources of low-grade lignite, a resource that is otherwise difficult to transport for its weight, ash content, low calorific value and propensity to spontaneous combustion during storage and transport. A number of similar plants exist around the world, although mostly these plants are targeted towards the production of long chain alkanes for use as gasoline, diesel, or feedstock to other processes.
Power to methane is a technology that uses electrical power to produce hydrogen from water by electrolysis and uses the Sabatier reaction to combine hydrogen with carbon dioxide to produce methane.
Laboratory synthesis
Methane can be produced by protonation of methyl lithium or a methyl Grignard reagent such as methylmagnesium chloride. It can also be made from anhydrous sodium acetate and dry sodium hydroxide, mixed and heated above 300 °C (with sodium carbonate as byproduct). In practice, a requirement for pure methane can easily be fulfilled by steel gas bottle from standard gas suppliers.
Occurrence
Methane is the major component of natural gas, about 87% by volume. The major source of methane is extraction from geological deposits known as natural gas fields, with coal seam gas extraction becoming a major source (see coal bed methane extraction, a method for extracting methane from a coal deposit, while enhanced coal bed methane recovery is a method of recovering methane from non-mineable coal seams). It is associated with other hydrocarbon fuels, and sometimes accompanied by helium and nitrogen. Methane is produced at shallow levels (low pressure) by anaerobic decay of organic matter and reworked methane from deep under the Earth's surface. In general, the sediments that generate natural gas are buried deeper and at higher temperatures than those that contain oil.
Methane is generally transported in bulk by pipeline in its natural gas form, or by LNG carriers in its liquefied form; few countries transport it by truck.
Atmospheric methane and climate change
Methane is an important greenhouse gas, responsible for around 30% of the rise in global temperatures since the industrial revolution.
Methane has a global warming potential (GWP) of 29.8 ± 11 compared to (potential of 1) over a 100-year period, and 82.5 ± 25.8 over a 20-year period. This means that, for example, a leak of one tonne of methane is equivalent to emitting 82.5 tonnes of carbon dioxide. Burning methane and producing carbon dioxide also reduces the greenhouse gas impact compared to simply venting methane to the atmosphere.
As methane is gradually converted into carbon dioxide (and water) in the atmosphere, these values include the climate forcing from the carbon dioxide produced from methane over these timescales.
Annual global methane emissions are currently approximately 580 Mt, 40% of which is from natural sources and the remaining 60% originating from human activity, known as anthropogenic emissions. The largest anthropogenic source is agriculture, responsible for around one quarter of emissions, closely followed by the energy sector, which includes emissions from coal, oil, natural gas and biofuels.
Historic methane concentrations in the world's atmosphere have ranged between 300 and 400 nmol/mol during glacial periods commonly known as ice ages, and between 600 and 700 nmol/mol during the warm interglacial periods. A 2012 NASA website said the oceans were a potential important source of Arctic methane, but more recent studies associate increasing methane levels as caused by human activity.
Global monitoring of atmospheric methane concentrations began in the 1980s. The Earth's atmospheric methane concentration has increased 160% since preindustrial levels in the mid-18th century. In 2013, atmospheric methane accounted for 20% of the total radiative forcing from all of the long-lived and globally mixed greenhouse gases. Between 2011 and 2019 the annual average increase of methane in the atmosphere was 1866 ppb. From 2015 to 2019 sharp rises in levels of atmospheric methane were recorded.
In 2019, the atmospheric methane concentration was higher than at any time in the last 800,000 years. As stated in the AR6 of the IPCC, "Since 1750, increases in (47%) and (156%) concentrations far exceed, and increases in (23%) are similar to, the natural multi-millennial changes between glacial and interglacial periods over at least the past 800,000 years (very high confidence)".
In February 2020, it was reported that fugitive emissions and gas venting from the fossil fuel industry may have been significantly underestimated.
The largest annual increase occurred in 2021 with the overwhelming percentage caused by human activity.
Climate change can increase atmospheric methane levels by increasing methane production in natural ecosystems, forming a climate change feedback. Another explanation for the rise in methane emissions could be a slowdown of the chemical reaction that removes methane from the atmosphere.
Over 100 countries have signed the Global Methane Pledge, launched in 2021, promising to cut their methane emissions by 30% by 2030. This could avoid 0.2˚C of warming globally by 2050, although there have been calls for higher commitments in order to reach this target. The International Energy Agency's 2022 report states "the most cost-effective opportunities for methane abatement are in the energy sector, especially in oil and gas operations".
Clathrates
Methane clathrates (also known as methane hydrates) are solid cages of water molecules that trap single molecules of methane. Significant reservoirs of methane clathrates have been found in arctic permafrost and along continental margins beneath the ocean floor within the gas clathrate stability zone, located at high pressures (1 to 100 MPa; lower end requires lower temperature) and low temperatures (< 15 °C; upper end requires higher pressure). Methane clathrates can form from biogenic methane, thermogenic methane, or a mix of the two. These deposits are both a potential source of methane fuel as well as a potential contributor to global warming. The global mass of carbon stored in gas clathrates is still uncertain and has been estimated as high as 12,500 Gt carbon and as low as 500 Gt carbon. The estimate has declined over time with a most recent estimate of ≈1800 Gt carbon. A large part of this uncertainty is due to our knowledge gap in sources and sinks of methane and the distribution of methane clathrates at the global scale. For example, a source of methane was discovered relatively recently in an ultraslow spreading ridge in the Arctic. Some climate models suggest that today's methane emission regime from the ocean floor is potentially similar to that during the period of the Paleocene–Eocene Thermal Maximum (PETM) around 55.5 million years ago, although there are no data indicating that methane from clathrate dissociation currently reaches the atmosphere. Arctic methane release from permafrost and seafloor methane clathrates is a potential consequence and further cause of global warming; this is known as the clathrate gun hypothesis. Data from 2016 indicate that Arctic permafrost thaws faster than predicted.
Public safety and the environment
Methane "degrades air quality and adversely impacts human health, agricultural yields, and ecosystem productivity".
Methane is extremely flammable and may form explosive mixtures with air. Methane gas explosions are responsible for many deadly mining disasters. A methane gas explosion was the cause of the Upper Big Branch coal mine disaster in West Virginia on April 5, 2010, killing 29. Natural gas accidental release has also been a major focus in the field of safety engineering, due to past accidental releases that concluded in the formation of jet fire disasters.
The 2015–2016 methane gas leak in Aliso Canyon, California was considered to be the worst in terms of its environmental effect in American history. It was also described as more damaging to the environment than Deepwater Horizon's leak in the Gulf of Mexico.
In May 2023 The Guardian published a report blaming Turkmenistan as the worst in the world for methane super emitting. The data collected by Kayrros researchers indicate that two large Turkmen fossil fuel fields leaked 2.6 million and 1.8 million metric tonnes of methane in 2022 alone, pumping the equivalent of 366 million tonnes into the atmosphere, surpassing the annual emissions of the United Kingdom.
Methane is also an asphyxiant if the oxygen concentration is reduced to below about 16% by displacement, as most people can tolerate a reduction from 21% to 16% without ill effects. The concentration of methane at which asphyxiation risk becomes significant is much higher than the 5–15% concentration in a flammable or explosive mixture. Methane off-gas can penetrate the interiors of buildings near landfills and expose occupants to significant levels of methane. Some buildings have specially engineered recovery systems below their basements to actively capture this gas and vent it away from the building.
Extraterrestrial methane
Interstellar medium
Methane is abundant in many parts of the Solar System and potentially could be harvested on the surface of another Solar System body (in particular, using methane production from local materials found on Mars or Titan), providing fuel for a return journey.
Mars
Methane has been detected on all planets of the Solar System and most of the larger moons. With the possible exception of Mars, it is believed to have come from abiotic processes.
The Curiosity rover has documented seasonal fluctuations of atmospheric methane levels on Mars. These fluctuations peaked at the end of the Martian summer at 0.6 parts per billion.
Methane has been proposed as a possible rocket propellant on future Mars missions due in part to the possibility of synthesizing it on the planet by in situ resource utilization. An adaptation of the Sabatier methanation reaction may be used with a mixed catalyst bed and a reverse water-gas shift in a single reactor to produce methane and oxygen from the raw materials available on Mars, utilizing water from the Martian subsoil and carbon dioxide in the Martian atmosphere.
Methane could be produced by a non-biological process called serpentinization involving water, carbon dioxide, and the mineral olivine, which is known to be common on Mars.
Titan
Methane has been detected in vast abundance on Titan, the largest moon of Saturn. It comprises a significant portion of its atmosphere and also exists in a liquid form on its surface, where it comprises the majority of the liquid in Titan's vast lakes of hydrocarbons, the second largest of which is believed to be almost pure methane in composition.
The presence of stable lakes of liquid methane on Titan, as well as the surface of Titan being highly chemically active and rich in organic compounds, has led scientists to consider the possibility of life existing within Titan's lakes, using methane as a solvent in the place of water for Earth-based life and using hydrogen in the atmosphere to derive energy with acetylene, in much the same way that Earth-based life uses glucose.
History
The discovery of methane is credited to Italian physicist Alessandro Volta, who characterized numerous properties including its flammability limit and origin from decaying organic matter.
Volta was initially motivated by reports of inflammable air present in marshes by his friend Father Carlo Guiseppe Campi. While on a fishing trip to Lake Maggiore straddling Italy and Switzerland in November 1776, he noticed the presence of bubbles in the nearby marshes and decided to investigate. Volta collected the gas rising from the marsh and demonstrated that the gas was inflammable.
Volta notes similar observations of inflammable air were present previously in scientific literature, including a letter written by Benjamin Franklin.
Following the Felling mine disaster of 1812 in which 92 men perished, Sir Humphry Davy established that the feared firedamp was in fact largely methane.
The name "methane" was coined in 1866 by the German chemist August Wilhelm von Hofmann. The name was derived from methanol.
Etymology
Etymologically, the word methane is coined from the chemical suffix "-ane", which denotes substances belonging to the alkane family; and the word methyl, which is derived from the German (1840) or directly from the French , which is a back-formation from the French (corresponding to English "methylene"), the root of which was coined by Jean-Baptiste Dumas and Eugène Péligot in 1834 from the Greek (wine) (related to English "mead") and (meaning "wood"). The radical is named after this because it was first detected in methanol, an alcohol first isolated by distillation of wood. The chemical suffix -ane is from the coordinating chemical suffix -ine which is from Latin feminine suffix -ina which is applied to represent abstracts. The coordination of "-ane", "-ene", "-one", etc. was proposed in 1866 by German chemist August Wilhelm von Hofmann.
Abbreviations
The abbreviation -C can mean the mass of carbon contained in a mass of methane, and the mass of methane is always 1.33 times the mass of -C. -C can also mean the methane-carbon ratio, which is 1.33 by mass. Methane at scales of the atmosphere is commonly measured in teragrams (Tg ) or millions of metric tons (MMT ), which mean the same thing. Other standard units are also used, such as nanomole (nmol, one billionth of a mole), mole (mol), kilogram, and gram.
| Physical sciences | Hydrocarbons | null |
18582511 | https://en.wikipedia.org/wiki/Alpaca | Alpaca | The alpaca (Lama pacos) is a species of South American camelid mammal. Traditionally, alpacas were kept in herds that grazed on the level heights of the Andes of Southern Peru, Western Bolivia, Ecuador, and Northern Chile. More recently, alpacas may be found on farms and ranches worldwide, with thousands of animals born and raised annually. Alpacas are especially popular in North America, Europe, and Australia.
There are two modern breeds of alpaca, separated based on their respective region of endemism and fiber (wool) type: the Suri alpaca and the Huacaya alpaca. Both breeds produce a highly valued fiber, with Suri alpaca's fiber growing in straight "locks," while Huacaya fiber has a "crimped," wavy texture and grows in bundles. These breeds' fibers are used for making knitted and woven items, similar to sheep's wool.
Alpacas are visually and genetically similar to, and often confused with a relative species, the llamas; however, alpacas are visibly shorter and predominantly bred for their wool, while llamas have long been more highly prized as livestock guardians (in place of dogs), and as a pack animal (beast-of-burden), owing to their nimble mountain-climbing abilities. Nonetheless, all four South American camelids are closely related and can successfully crossbreed. Both the alpaca and the llama are believed to have been domesticated and selectively bred from their wild counterparts — the smaller, fine-haired vicuña and the larger, stronger guanaco, respectively — at least 5,000 to 6,000 years ago.
Alpacas communicate through body language, spitting to show dominance when distressed, fearful, or agitated. Male alpacas are more aggressive than females. In some cases, alpha males will immobilize the head and neck of a weaker or challenging male to show their strength and dominance.
In the textile industry, "alpaca" primarily refers to the hair of Peruvian alpacas. More broadly, it refers to a style of fabric originally made from alpaca hair, such as mohair, Icelandic sheep wool, or even high-quality wool from other breeds of sheep. In trade, distinctions are made between alpacas and the several styles of mohair and luster.
Background
The relationship between alpacas and vicuñas was disputed for many years. In the 18th and 19th centuries, the four South American lamoid species were assigned scientific names. At that time, the alpaca was assumed to be descended from the llama, ignoring similarities in size, fleece, and dentition between the alpaca and the vicuña. Classification was complicated by the fact that all four species of South American camelid can interbreed and produce fertile offspring. The advent of DNA technology made a more accurate classification possible.
In 2001, the alpaca genus classification changed from Lama pacos to Vicugna pacos, following the presentation of a paper on work by Miranda Kadwell etal. on alpaca DNA to the Royal Society showing the alpaca is descended from the vicuña, not the guanaco.
An adult alpaca generally is between in height at the shoulders (withers). They usually weigh between . Raised in the same conditions, the difference in weight can be small with males weighing around and females .
Origin and domestication
Alpacas were domesticated thousands of years ago. The Moche people of Northern Peru often used alpaca images in their art. Traditionally, alpaca were bred and raised in herds, grazing on the level meadows and escarpments of the Andes, from Ecuador and Peru to Western Bolivia and Northern Chile, typically at an altitude of above sea level. There are no known wild alpacas, and its closest living relative, the vicuña (also native to South America), is the wild ancestor of the alpaca.
The family Camelid first appeared in the Americas 40–45 million years ago, during the Eocene period, from the common ancestor, Protylopus. The descendants divided into Camelini and Lamini tribes, taking different migratory patterns to Asia and South America, respectively. Although the camelids became extinct in North America around 3 million years ago, they flourished in the South. It was not until 2–5 million years ago, during the Pliocene, that the genus Hemiauchenia of the tribe Lamini split into Palaeolama and Lama; the latter would then split again into Lama and Vicugna upon migrating down to South America.
Remains of vicuña and guanaco dating around 12,000 years have been found throughout Peru. Their domesticated counterparts, the llama and alpaca, have been found mummified in the Moquegua valley, in the South of Peru, dating back 900 to 1000 years. Mummies found in this region show two breeds of alpacas. More precise analysis of bone and teeth of these mummies has demonstrated that alpacas were domesticated from the Vicugna vicugna. Other research, considering the behavioral and morphological characteristics of alpacas and their wild counterparts, seems to indicate that alpacas could find their origins in Lama guanicoe as well as Vicugna vicugna, or even a hybrid of both.
Genetic analysis shows a different picture of the origins of the alpaca. Analysis of mitochondrial DNA shows that most alpacas have guanaco mtDNA, and many also have vicuña mtDNA. But microsatellite data shows that alpaca DNA is much more similar to vicuña DNA than to guanaco DNA. This suggests that alpacas are descendants of the Vicugna vicugna, not of the Lama guanicoe. The discrepancy with mtDNA seems to be because mtDNA is only transmitted by the mother, and recent husbandry practices have caused hybridization between llamas (which primarily carry guanaco DNA) and alpacas. To the extent that many domestic alpacas are the result of male alpacas bred to female llamas, this would explain the mtDNA consistent with guanacos. This situation has led to attempts to reclassify the alpaca as Vicugna pacos.
Breeds
The alpaca is divided into two breeds, Suri and Huacaya, based on its fibers rather than scientific or European classifications.
Huacaya alpacas are the most commonly found, constituting about 90% of the population. The Huacaya alpaca is thought to have originated in post-colonial Peru because of its thicker fleece which makes them more suited to survive in the higher altitudes of the Andes after being pushed into the highlands of Peru with the arrival of the Spanish.
Suri alpacas represent a smaller portion of the total alpaca population, around 10%. They are thought to have been more prevalent in pre-Columbian Peru since they could be kept at a lower altitude where a thicker fleece was not needed for harsh weather conditions.
Behavior
Alpacas are social herd animals that live in family groups, consisting of a territorial alpha male, females, and their young ones. Alpacas warn the herd about intruders by making sharp, noisy inhalations that sound like a high-pitched bray. The herd may attack smaller predators with their front feet and can spit and kick. Their aggression towards members of the canid family (coyotes, foxes, dogs, etc.) is exploited when alpacas are used as guard llamas for guarding sheep.
Alpacas can sometimes be aggressive, but they can also be very gentle, intelligent, and extremely observant. For the most part, alpacas are very quiet, but male alpacas are more energetic when they get involved in fighting with other alpacas. When they prey, they are cautious and nervous when they feel threatened. They can feel threatened when someone or another alpaca comes up behind them.
Alpacas set their boundaries of "personal space" within their families and groups. They make a hierarchy in some sense, and each alpaca is aware of the dominant animals in each group. Body language is the key to their communication. It helps to maintain their order. One example of their body communication includes a pose named broadside, where their ears are pulled back and they stand sideways. This pose is used when male alpacas are defending their territory. They commonly spit to show dominance when they are in distress, fearful, or feel agitated. Male alpacas are more aggressive than females and tend to establish dominance within their herd group. In some cases, alpha males will immobilize the head and neck of a weaker or challenging male to show their strength and dominance.
When they are young, alpacas tend to follow larger objects and sit near or under them. An example of this is a baby alpaca with its mother. This can also apply when an alpaca passes by an older alpaca.
Training
Alpacas are generally very trainable and usually respond to rewards, most commonly food. They can usually be petted without getting agitated, especially if one avoids petting the head or neck. Alpacas are usually quite easy to herd, even in large groups. However, during herding, it is recommended for the handler to approach the animals slowly and quietly, as failing to do so can result in danger for both the animals and the handler.
Alpacas and llamas have started showing up in U.S. nursing homes and hospitals as trained, certified therapy animals. The Mayo Clinic says animal-assisted therapy can reduce pain, depression, anxiety, and fatigue. This type of animal therapy is growing in popularity, and several organizations throughout the United States participate.
Spitting
Not all alpacas spit, but all are capable of doing so. "Spit" is somewhat euphemistic; occasionally, the projectile contains only air and a little saliva, although alpacas commonly bring up acidic stomach contents (generally a green, grassy mix) and project it onto their chosen targets. Spitting is mostly reserved for other alpacas, but an alpaca will also occasionally spit at a human.
Spitting can result in what is called a "sour mouth." A sour mouth is characterized by "a loose-hanging lower lip and a gaping mouth."
Alpacas can spit for several reasons. A female alpaca spits when she is not interested in a male alpaca, typically when she thinks that she is already impregnated. Both sexes of alpaca keep each other away from their food or anything they have their eyes on. Most give a slight warning before spitting by blowing air out and raising their heads, giving their ears a "pinned" appearance.
Alpacas can spit up to ten feet if they need to. For example, if another animal does not back off, the alpaca will throw up its stomach contents, resulting in a lot of spit.
Some signs of stress that can lead to their spitting habits include humming, a wrinkle under their eye, drooling, rapid breathing, and stomping their feet. When alpacas show any sign of interest or alertness, they tend to sniff their surroundings, watch closely, or stand quietly in place and stare.
When it comes to reproduction, they spit because it is a response triggered by the increased progesterone levels, which is associated with ovulation.
Hygiene
Alpacas use a communal dung pile, where they do not graze. This behavior may limit the spread of internal parasites. Generally, males have much tidier and fewer dung piles than females, which tend to stand in a line and all go simultaneously. One female approaches the dung pile and begins to urinate and/or defecate, and the rest of the herd often follows. Alpaca waste is collected and can be used as soil fertilizer.
Because they prefer using a dung pile for excreting bodily waste, some alpacas have been successfully house-trained.
Alpacas develop dental hygiene problems that affect their eating and behavior. Warning signs include protracted chewing while eating or food spilling out of their mouths. Poor body condition and sunken cheeks are also telltale signs of dental problems.
Sounds
Alpacas make a variety of sounds:
Humming: When alpacas are born, the mother and baby hum constantly. They also hum as a sign of distress, especially when they are separated from their herd. Alpacas may also hum when curious, happy, worried, or cautious.
Snorting: Alpacas snort when another alpaca is invading their space.
Grumbling: Alpacas grumble to warn each other. For example, when one is invading another's personal space, it sounds like gurgling.
Clucking: Similar to a hen's cluck, alpacas cluck when a mother is concerned for her cria. Male alpacas cluck to signal friendly behavior.
Screaming: Their screams are extremely deafening and loud. They will scream when they are not handled correctly or when a potential enemy is attacking them.
Screeching: A bird-like cry, presumably intended to terrify the opponent. This sound is typically used by male alpacas when they are in a fight over dominance. When a female screeches, it is more of a growl when she is angry.
Reproduction
Females are induced ovulators; meaning that the act of mating and the presence of semen causes them to ovulate. Females usually conceive after just one breeding but occasionally do have trouble conceiving. Artificial insemination is technically difficult, expensive, and uncommon but feasible. On the contrary, embryo transfer is more widespread.
A male is usually ready to mate for the first time between two and three years of age. It is not advisable to allow a young female to be bred until she is mature and has reached two-thirds of her mature weight. Overbreeding a young female before conception is possibly a common cause of uterine infections. As the age of maturation varies greatly between individuals, it is usually recommended that novice breeders wait until females are 18 months of age or older before initiating breeding.
Alpacas can breed at any time throughout the year, but it is more difficult to breed in the winter. Most breed during autumn or late spring. The most popular way to have alpacas mate is pen mating, which involves moving both the female and the desired male into a pen. Another way is paddock mating, where one male alpaca is let loose in the paddock with several female alpacas.
The gestation period is, on average, 11.5 months, and usually results in a single offspring, or cria. Twins are rare, occurring about once per 1,000 deliveries. Cria are generally between 15 and 19 pounds, and are standing 30 to 90 minutes after birth. Two weeks after a female gives birth, she is generally receptive to breeding again. Crias may be weaned through human intervention at about six months old and 60 pounds weight. However, many breeders prefer to allow the female to decide when to wean her offspring; they can be weaned earlier or later depending on their size and emotional maturity.
The average lifespan of an alpaca is between 15 and 20 years, and the longest-lived alpaca on record is 28 years.
Pests and diseases
Cattle tuberculosis can also infect alpacas: Mycobacterium bovis also causes TB in this species worldwide. Krajewska-Wędzina et al., 2020 detect M. bovis in individuals traded from the United Kingdom to Poland. To accomplish this they develop a seroassay which correctly identifies positive subjects which are false negative for a common skin test. Krajewska-Wędzina et al. also find that alpacas are unusual in mounting a competent early-infection immune response. Bernitz et al., 2021 believe this to generalise to all camelids.
Habitat and lifestyle
Alpacas can be found throughout most of South America. They typically live in temperate conditions in the mountains with high altitudes.
They are easy to care for since they are not limited to a specific type of environment. Animals such as flamingos, condors, spectacled bears, mountain lions, coyotes, llamas, and sheep live near alpacas when they are in their natural habitat.
Population
Alpacas are native to Peru but can be found throughout the globe in captivity. Peru currently has the largest alpaca population, with over half the world's animals. The population declined drastically after the Spanish Conquistadors invaded the Andes mountains in 1532, after which 98% of the animals were destroyed. The Spanish also brought with them diseases that were fatal to alpacas.
European conquest forced the animals to move higher into the mountains, which remained there permanently. Although alpacas had almost been wiped out completely, they were rediscovered sometime during the 19th century by Europeans. After finding their uses, animals became important to societies during the Industrial Revolution.
Diet
Alpacas chew their food which ends up being mixed with their cud and saliva and then they swallow it. Alpacas usually eat 1.5% of their body weight daily for normal growth. They mainly need pasture grass, hay, or silage. Still, some may also need supplemental energy and protein foods, and they will also usually try to chew on almost anything (e.g., empty bottles). Most alpaca ranchers rotate their feeding grounds so the grass can regrow, and fecal parasites may die before reusing the area. Pasture grass is a great source of protein. When seasons change, the grass loses or gains more protein. For example, in the spring, the pasture grass has about 20% protein, while it only has 6% in the summer. They need more energy supplements in the winter to produce body heat and warmth. They get their fiber from hay or long stems, which provides them with vitamin E. Green grass contains vitamins A and E.
Alpacas can eat natural unfertilized grass; however, ranchers can also supplement grass with low-protein grass hay. To provide selenium and other necessary vitamins, ranchers will feed their domestic alpacas a daily dose of grain to provide additional nutrients that are not fully obtained from their primary diet. Alpacas may obtain the necessary vitamins in their native grazing ranges.
Digestion
Like other camelids, alpacas have a three-chambered stomach; combined with chewing cud, this three-chambered system allows maximum extraction of nutrients from low-quality forages. Alpacas are not ruminants, pseudo-ruminants, or modified ruminants, as there are many differences between the anatomy and physiology of a camelid and a ruminant stomach.
Alpacas chew their food in a figure eight motion, swallow it, and then pass it into one of their stomach's chambers. The first and second chambers (C1 and C2) are anaerobic fermentation chambers where the fermentation process begins. The alpaca will further absorb nutrients and water in the first part of the third chamber. The end of the third chamber (called C3) is where the stomach secretes acids to digest food and is the likely place where an alpaca will have ulcers if stressed.
Poisonous plants
Unlike sheep and goats, which are commonly used to clear overgrown patches of land—as they willingly consume many noxious, poisonous botanical species—, many more common plant families are highly poisonous to alpacas, including the Amaryllidaceae (Amaryllis) family, the dogbane-milkweed family Apocynaceae (Asclepias, Hoya, Nerium, Plumeria, etc), the aroid family Araceae (Anthurium, Colocasia, Monstera, Philodendron, Zantedeschia, etc), the Asparagaceae (Agave, Asparagus, Dracaena, and more), Asteraceae (daisies and Senecio, etc), Caryophyllaceae (Dianthus), some Ericaceae (azaleas, heather, etc), Euphorbiaceae (castor bean, Croton, poinsettia, etc), Fagaceae (beech and oak; acorns), ferns (especially Pteridium), African rue, Iridaceae (Crocus, Freesia, Gladiolus, Iris, etc), Melanthiaceae (corn-lilies), Polygonaceae (buckwheat, knotweed), ragweed, Ranunculaceae (buttercups), as well as orange tree and other Citrus foliage, among others.
Fiber
Alpaca fleece is soft and possesses water and flame resistant properties, making it a valuable commodity. It is used for making knitted and woven items, similar to sheep's wool. These items include blankets, sweaters, hats, gloves, scarves, a wide variety of textiles, and ponchos, in South America, as well as sweaters, socks, coats, and bedding in other parts of the world. The fiber comes in more than 52 natural colors as classified in Peru, 12 as classified in Australia, and 16 as classified in the United States.
Alpacas are typically sheared once per year in the spring. Each shearing produces approximately of fiber per alpaca. An adult alpaca might produce of first-quality fiber as well as of second- and third-quality fiber. The quality of alpaca fiber is determined by how crimpy it is. Typically, the greater the number of small folds in the fiber, the greater the quality.
There are two modern breeds of alpaca, separated based on their respective region of endemism and fiber (wool) type—the Suri alpaca and the Huacaya alpaca. Both breeds produce a fiber that is highly valued, with Suri alpaca fiber growing in straight "locks" and being comparable to the mohair of Turkish Angora goats; Huacaya has a "crimped", wavy texture and grows in bundles more similar to sheep's wool.
Prices
Alpacas were the subject of a speculative bubble between their introduction to North America in 1984 and the early 21st century. The price for American alpacas ranged from US$50 for a castrated male (gelding) to US$675,000 for the highest in the world, depending on breeding history, sex, and color. In 2006, researchers warned that the higher prices sought for alpaca breeding stock were largely speculative and not supported by market fundamentals, given the low inherent returns per head from the main end product, alpaca fiber, and prices into the $100s per head rather than $10,000s would be required for a commercially viable fiber production herd.
Marketed as "the investment you can hug" in television commercials by the Alpaca Owners and Breeders Association, the market for alpacas was almost entirely dependent on breeding and selling animals to new buyers, a classic sign of speculative bubbles in livestock. The bubble burst in 2007, with the price of alpaca breeding stock dropping by thousands of dollars each year thereafter. Many farmers found themselves unable to sell animals for any price, or even give them away.
It is possible to raise up to , as they have a designated area for waste products and keep their eating area away from their waste area. However, this ratio differs from country to country and is highly dependent on the quality of pasture available (in many desert locations it is generally only possible to run one to three animals per acre due to lack of suitable vegetation). Fiber quality is the primary variant in the price achieved for alpaca wool; in Australia, it is common to classify the fiber by the thickness of the individual hairs and by the amount of vegetable matter contained in the supplied shearings.
Livestock
Alpacas need to eat 1–2% of their body weight per day, so about two bales of grass hay per month per animal. When formulating a proper diet for alpacas, water and hay analysis should be performed to determine the proper vitamin and mineral supplementation program. Two options are to provide free choice salt/mineral powder or feed a specially formulated ration. Indigenous to the highest regions of the Andes, this harsh environment has created an extremely hardy animal, so only minimal housing and predator fencing are needed. The alpacas' three-chambered stomachs allow for extremely efficient digestion. There are no viable seeds in the manure, because alpacas prefer to only eat tender plant leaves, and will not consume thick plant stems; therefore, alpaca manure does not need composting to enrich pastures or ornamental landscaping. Nail and teeth trimming are needed every six to twelve months, along with annual shearing.
Similar to ruminants, such as cattle and sheep, alpacas have only lower teeth at the front of their mouths; therefore, they do not pull the grass up by the roots. Rotating pastures is still important, though, as alpacas have a tendency to regraze an area repeatedly. Alpacas are fiber-producing animals; they do not need to be slaughtered to reap their product, and their fiber is a renewable resource that grows yearly.
Cultural presence
Alpacas are closely tied to cultural practices for Andeans people. Prior to colonization, the image of the alpaca was used in rituals and in their religious practices. Since the people in the region depended heavily on these animals for their sustenance, the alpaca was seen as a gift from Pachamama. Alpacas were used for their meat, fibers for clothing, and art, and their images in the form of conopas.
Conopas take their appearance from the Suri alpacas, with long locks flanking their sides and bangs covering the eyes, and a depression on the back. This depression is used in ritual practices, usually filled with coca leaves and fat from alpacas and lamas, to bring fertility and luck. While their use was prevalent before colonization, the attempts to convert the Andean people to Catholicism led to the acquisition of more than 3,400 conopas in Lima alone.
The origin of alpacas is depicted in legend; the legend states they came to be in the world after a goddess fell in love with a man. The goddess' father only allowed her to be with her lover if he cared for her herd of alpacas. On top of caring for the herd, he was to always carry a small animal for his entire life. As the goddess came into our world, the alpacas followed her. Everything was fine until the man set the small animal down, and the goddess fled back to her home. On her way back home, the man attempted to stop her and her herd from fleeing. While he was not able to stop her from returning, he was able to stop a few alpacas from returning. These alpacas who did not make it back are said to be seen today in the swampy lands in the Andes waiting for the end of the world, so they may return to their goddess.
| Biology and health sciences | Artiodactyla | null |
18588994 | https://en.wikipedia.org/wiki/Usenet | Usenet | Usenet (), USENET, or, "in full", User's Network, is a worldwide distributed discussion system available on computers. It was developed from the general-purpose Unix-to-Unix Copy (UUCP) dial-up network architecture. Tom Truscott and Jim Ellis conceived the idea in 1979, and it was established in 1980. Users read and post messages (called articles or posts, and collectively termed news) to one or more topic categories, known as newsgroups. Usenet resembles a bulletin board system (BBS) in many respects and is the precursor to the Internet forums that have become widely used. Discussions are threaded, as with web forums and BBSes, though posts are stored on the server sequentially.
A major difference between a BBS or web message board and Usenet is the absence of a central server and dedicated administrator or hosting provider. Usenet is distributed among a large, constantly changing set of news servers that store and forward messages to one another via "news feeds". Individual users may read messages from and post to a local (or simply preferred) news server, which can be operated by anyone, and those posts will automatically be forwarded to any other news servers peered with the local one, while the local server will receive any news its peers have that it currently lacks. This results in the automatic proliferation of content posted by any user on any server to any other user subscribed to the same newsgroups on other servers.
As with BBSes and message boards, individual news servers or service providers are under no obligation to carry any specific content, and may refuse to do so for many reasons: a news server might attempt to control the spread of spam by refusing to accept or forward any posts that trigger spam filters, or a server without high-capacity data storage may refuse to carry any newsgroups used primarily for file sharing, limiting itself to discussion-oriented groups. However, unlike BBSes and web forums, the dispersed nature of Usenet usually permits users who are interested in receiving some content to access it simply by choosing to connect to news servers that carry the feeds they want.
Usenet is culturally and historically significant in the networked world, having given rise to, or popularized, many widely recognized concepts and terms such as "FAQ", "flame", "sockpuppet", and "spam". In the early 1990s, shortly before access to the Internet became commonly affordable, Usenet connections via FidoNet's dial-up BBS networks made long-distance or worldwide discussions and other communication widespread, not needing a server, just (local) telephone service.
The name Usenet comes from the term "users' network". The first Usenet group was NET.general, which quickly became net.general. The first commercial spam on Usenet was from immigration attorneys Canter and Siegel advertising green card services.
On the Internet, Usenet is transported via the Network News Transfer Protocol (NNTP) on Transmission Control Protocol (TCP) port 119 for standard, unprotected connections, and on TCP port 563 for Secure Sockets Layer (SSL) encrypted connections.
Introduction
Usenet was conceived in 1979 and publicly established in 1980, at the University of North Carolina at Chapel Hill and Duke University, over a decade before the World Wide Web went online (and thus before the general public received access to the Internet), making it one of the oldest computer network communications systems still in widespread use. It was originally built on the "poor man's ARPANET", employing UUCP as its transport protocol to offer mail and file transfers, as well as announcements through the newly developed news software such as A News. The name "Usenet" emphasizes its creators' hope that the USENIX organization would take an active role in its operation.
The articles that users post to Usenet are organized into topical categories known as newsgroups, which are themselves logically organized into hierarchies of subjects. For instance, sci.math and sci.physics are within the sci.* hierarchy. Or, talk.origins and talk.atheism are in the talk.* hierarchy. When a user subscribes to a newsgroup, the news client software keeps track of which articles that user has read.
In most newsgroups, the majority of the articles are responses to some other article. The set of articles that can be traced to one single non-reply article is called a thread. Most modern newsreaders display the articles arranged into threads and subthreads. For example, in the wine-making newsgroup rec.crafts.winemaking, someone might start a thread called; "What's the best yeast?" and that thread or conversation might grow into dozens of replies long, by perhaps six or eight different authors. Over several days, that conversation about different wine yeasts might branch into several sub-threads in a tree-like form.
When a user posts an article, it is initially only available on that user's news server. Each news server talks to one or more other servers (its "newsfeeds") and exchanges articles with them. In this fashion, the article is copied from server to server and should eventually reach every server in the network. The later peer-to-peer networks operate on a similar principle, but for Usenet it is normally the sender, rather than the receiver, who initiates transfers. Usenet was designed under conditions when networks were much slower and not always available. Many sites on the original Usenet network would connect only once or twice a day to batch-transfer messages in and out. This is largely because the POTS network was typically used for transfers, and phone charges were lower at night.
The format and transmission of Usenet articles is similar to that of Internet e-mail messages. The difference between the two is that Usenet articles can be read by any user whose news server carries the group to which the message was posted, as opposed to email messages, which have one or more specific recipients.
Today, Usenet has diminished in importance with respect to Internet forums, blogs, mailing lists and social media. Usenet differs from such media in several ways: Usenet requires no personal registration with the group concerned; information need not be stored on a remote server; archives are always available; and reading the messages does not require a mail or web client, but a news client. However, it is now possible to read and participate in Usenet newsgroups to a large degree using ordinary web browsers since most newsgroups are now copied to several web sites. The groups in are still widely used for data transfer.
ISPs, news servers, and newsfeeds
Many Internet service providers, and many other Internet sites, operate news servers for their users to access. ISPs that do not operate their own servers directly will often offer their users an account from another provider that specifically operates newsfeeds. In early news implementations, the server and newsreader were a single program suite, running on the same system. Today, one uses separate newsreader client software, a program that resembles an email client but accesses Usenet servers instead.
Not all ISPs run news servers. A news server is one of the most difficult Internet services to administer because of the large amount of data involved, small customer base (compared to mainstream Internet service), and a disproportionately high volume of customer support incidents (frequently complaining of missing news articles). Some ISPs outsource news operations to specialist sites, which will usually appear to a user as though the ISP itself runs the server. Many of these sites carry a restricted newsfeed, with a limited number of newsgroups. Commonly omitted from such a newsfeed are foreign-language newsgroups and the hierarchy which largely carries software, music, videos and images, and accounts for over 99 percent of article data.
There are also Usenet providers that offer a full unrestricted service to users whose ISPs do not carry news, or that carry a restricted feed.
Newsreaders
Newsgroups are typically accessed with newsreaders: applications that allow users to read and reply to postings in newsgroups. These applications act as clients to one or more news servers. Historically, Usenet was associated with the Unix operating system developed at AT&T, but newsreaders were soon available for all major operating systems. Email client programs and Internet suites of the late 1990s and 2000s often included an integrated newsreader. Newsgroup enthusiasts often criticized these as inferior to standalone newsreaders that made correct use of Usenet protocols, standards and conventions.
With the rise of the World Wide Web (WWW), web front-ends (web2news) have become more common. Web front ends have lowered the technical entry barrier requirements to that of one application and no Usenet NNTP server account. There are numerous websites now offering web based gateways to Usenet groups, although some people have begun filtering messages made by some of the web interfaces for one reason or another. Google Groups is one such web based front end and some web browsers can access Google Groups via news: protocol links directly.
Moderated and unmoderated newsgroups
A minority of newsgroups are moderated, meaning that messages submitted by readers are not distributed directly to Usenet, but instead are emailed to the moderators of the newsgroup for approval. The moderator is to receive submitted articles, review them, and inject approved articles so that they can be properly propagated worldwide. Articles approved by a moderator must bear the Approved: header line. Moderators ensure that the messages that readers see in the newsgroup conform to the charter of the newsgroup, though they are not required to follow any such rules or guidelines. Typically, moderators are appointed in the proposal for the newsgroup, and changes of moderators follow a succession plan.
Historically, a mod.* hierarchy existed before Usenet reorganization. Now, moderated newsgroups may appear in any hierarchy, typically with .moderated added to the group name.
Usenet newsgroups in the Big-8 hierarchy are created by proposals called a Request for Discussion, or RFD. The RFD is required to have the following information: newsgroup name, checkgroups file entry, and moderated or unmoderated status. If the group is to be moderated, then at least one moderator with a valid email address must be provided. Other information which is beneficial but not required includes: a charter, a rationale, and a moderation policy if the group is to be moderated. Discussion of the new newsgroup proposal follows, and is finished with the members of the Big-8 Management Board making the decision, by vote, to either approve or disapprove the new newsgroup.
Unmoderated newsgroups form the majority of Usenet newsgroups, and messages submitted by readers for unmoderated newsgroups are immediately propagated for everyone to see. Minimal editorial content filtering vs propagation speed form one crux of the Usenet community. One little cited defense of propagation is canceling a propagated message, but few Usenet users use this command and some news readers do not offer cancellation commands, in part because article storage expires in relatively short order anyway. Almost all unmoderated Usenet groups tend to receive large amounts of spam.
Technical details
Usenet is a set of protocols for generating, storing and retrieving news "articles" (which resemble Internet mail messages) and for exchanging them among a readership which is potentially widely distributed. These protocols most commonly use a flooding algorithm which propagates copies throughout a network of participating servers. Whenever a message reaches a server, that server forwards the message to all its network neighbors that haven't yet seen the article. Only one copy of a message is stored per server, and each server makes it available on demand to the (typically local) readers able to access that server. The collection of Usenet servers has thus a certain peer-to-peer character in that they share resources by exchanging them, the granularity of exchange however is on a different scale than a modern peer-to-peer system and this characteristic excludes the actual users of the system who connect to the news servers with a typical client-server application, much like an email reader.
RFC 850 was the first formal specification of the messages exchanged by Usenet servers. It was superseded by RFC 1036 and subsequently by RFC 5536 and RFC 5537.
In cases where unsuitable content has been posted, Usenet has support for automated removal of a posting from the whole network by creating a cancel message, although due to a lack of authentication and resultant abuse, this capability is frequently disabled. Copyright holders may still request the manual deletion of infringing material using the provisions of World Intellectual Property Organization treaty implementations, such as the United States Online Copyright Infringement Liability Limitation Act, but this would require giving notice to each individual news server administrator.
On the Internet, Usenet is transported via the Network News Transfer Protocol (NNTP) on TCP Port 119 for standard, unprotected connections and on TCP port 563 for SSL encrypted connections.
Organization
The major set of worldwide newsgroups is contained within nine hierarchies, eight of which are operated under consensual guidelines that govern their administration and naming. The current Big Eight are:
comp.* – computer-related discussions (comp.software, comp.sys.amiga)
humanities.* – fine arts, literature, and philosophy (humanities.classics, humanities.design.misc)
misc.* – miscellaneous topics (misc.education, misc.forsale, misc.kids)
news.* – discussions and announcements about news (meaning Usenet, not current events) (news.groups, news.admin)
rec.* – recreation and entertainment (rec.music, rec.arts.movies)
sci.* – science related discussions (sci.psychology, sci.research)
soc.* – social discussions (soc.college.org, soc.culture.)
talk.* – talk about various controversial topics (talk.religion, talk.politics, talk.origins)
The alt.* hierarchy is not subject to the procedures controlling groups in the Big Eight, and it is as a result less organized. Groups in the alt.* hierarchy tend to be more specialized or specific—for example, there might be a newsgroup under the Big Eight which contains discussions about children's books, but a group in the alt hierarchy may be dedicated to one specific author of children's books. Binaries are posted in , making it the largest of all the hierarchies.
Many other hierarchies of newsgroups are distributed alongside these. Regional and language-specific hierarchies such as .*, .* and ne.* serve specific countries and regions such as Japan, Malta and New England. Companies and projects administer their own hierarchies to discuss their products and offer community technical support, such as the historical .* hierarchy from the Free Software Foundation. Microsoft closed its newsserver in June 2010, providing support for its products over forums now. Some users prefer to use the term "Usenet" to refer only to the Big Eight hierarchies; others include alt.* as well. The more general term "netnews" incorporates the entire medium, including private organizational news systems.
Informal sub-hierarchy conventions also exist. *.answers are typically moderated cross-post groups for FAQs. An FAQ would be posted within one group and a cross post to the *.answers group at the head of the hierarchy seen by some as a refining of information in that news group. Some subgroups are recursive—to the point of some silliness in alt.*.
Binary content
Usenet was originally created to distribute text content encoded in the 7-bit ASCII character set. With the help of programs that encode 8-bit values into ASCII, it became practical to distribute binary files as content. Binary posts, due to their size and often-dubious copyright status, were in time restricted to specific newsgroups, making it easier for administrators to allow or disallow the traffic.
The oldest widely used encoding method for binary content is uuencode, from the Unix UUCP package. In the late 1980s, Usenet articles were often limited to 60,000 characters, and larger hard limits exist today. Files are therefore commonly split into sections that require reassembly by the reader.
With the header extensions and the Base64 and Quoted-Printable MIME encodings, there was a new generation of binary transport. In practice, MIME has seen increased adoption in text messages, but it is avoided for most binary attachments. Some operating systems with metadata attached to files use specialized encoding formats. For Mac OS, both BinHex and special MIME types are used. Other lesser known encoding systems that may have been used at one time were BTOA, XX encoding, BOO, and USR encoding.
In an attempt to reduce file transfer times, an informal file encoding known as yEnc was introduced in 2001. It achieves about a 30% reduction in data transferred by assuming that most 8-bit characters can safely be transferred across the network without first encoding into the 7-bit ASCII space. The most common method of uploading large binary posts to Usenet is to convert the files into RAR archives and create Parchive files for them. Parity files are used to recreate missing data when not every part of the files reaches a server.
Binary newsgroups can be used to distribute files, and, as of 2022, some remain popular as an alternative to BitTorrent to share and download files.
Binary retention time
Each news server allocates a certain amount of storage space for content in each newsgroup. When this storage has been filled, each time a new post arrives, old posts are deleted to make room for the new content. If the network bandwidth available to a server is high but the storage allocation is small, it is possible for a huge flood of incoming content to overflow the allocation and push out everything that was in the group before it. The average length of time that posts are able to stay on the server before being deleted is commonly called the retention time.
Binary newsgroups are only able to function reliably if there is sufficient storage allocated to handle the amount of articles being added. Without sufficient retention time, a reader will be unable to download all parts of the binary before it is flushed out of the group's storage allocation. This was at one time how posting undesired content was countered; the newsgroup would be flooded with random garbage data posts, of sufficient quantity to push out all the content to be suppressed. This has been compensated by service providers allocating enough storage to retain everything posted each day, including spam floods, without deleting anything.
Modern Usenet news servers have enough capacity to archive years of binary content even when flooded with new data at the maximum daily speed available.
In part because of such long retention times, as well as growing Internet upload speeds, Usenet is also used by individual users to store backup data. While commercial providers offer easier to use online backup services, storing data on Usenet is free of charge (although access to Usenet itself may not be). The method requires the uploader to cede control over the distribution of the data; the files are automatically disseminated to all Usenet providers exchanging data for the news group it is posted to. In general the user must manually select, prepare and upload the data. The data is typically encrypted because it is available to anyone to download the backup files. After the files are uploaded, having multiple copies spread to different geographical regions around the world on different news servers decreases the chances of data loss.
Major Usenet service providers have a retention time of more than 12 years. This results in more than 60 petabytes (60000 terabytes) of storage (see image). When using Usenet for data storage, providers that offer longer retention time are preferred to ensure the data will survive for longer periods of time compared to services with lower retention time.
Legal issues
While binary newsgroups can be used to distribute completely legal user-created works, free software, and public domain material, some binary groups are used to illegally distribute proprietary software, copyrighted media, and pornographic material.
ISP-operated Usenet servers frequently block access to all groups to both reduce network traffic and to avoid related legal issues. Commercial Usenet service providers claim to operate as a telecommunications service, and assert that they are not responsible for the user-posted binary content transferred via their equipment. In the United States, Usenet providers can qualify for protection under the DMCA Safe Harbor regulations, provided that they establish a mechanism to comply with and respond to takedown notices from copyright holders.
Removal of copyrighted content from the entire Usenet network is a nearly impossible task, due to the rapid propagation between servers and the retention done by each server. Petitioning a Usenet provider for removal only removes it from that one server's retention cache, but not any others. It is possible for a special post cancellation message to be distributed to remove it from all servers, but many providers ignore cancel messages by standard policy, because they can be easily falsified and submitted by anyone. For a takedown petition to be most effective across the whole network, it would have to be issued to the origin server to which the content has been posted, before it has been propagated to other servers. Removal of the content at this early stage would prevent further propagation, but with modern high speed links, content can be propagated as fast as it arrives, allowing no time for content review and takedown issuance by copyright holders.
Establishing the identity of the person posting illegal content is equally difficult due to the trust-based design of the network. Like SMTP email, servers generally assume the header and origin information in a post is true and accurate. However, as in SMTP email, Usenet post headers are easily falsified so as to obscure the true identity and location of the message source. In this manner, Usenet is significantly different from modern P2P services; most P2P users distributing content are typically immediately identifiable to all other users by their network address, but the origin information for a Usenet posting can be completely obscured and unobtainable once it has propagated past the original server.
Also unlike modern P2P services, the identity of the downloaders is hidden from view. On P2P services a downloader is identifiable to all others by their network address. On Usenet, the downloader connects directly to a server, and only the server knows the address of who is connecting to it. Some Usenet providers do keep usage logs, but not all make this logged information casually available to outside parties such as the Recording Industry Association of America. The existence of anonymising gateways to USENET also complicates the tracing of a postings true origin.
History
UUCP/Usenet Logical Map, original by Steven McGeady.
Newsgroup experiments first occurred in 1979. Tom Truscott and Jim Ellis of Duke University came up with the idea as a replacement for a local announcement program, and established a link with nearby University of North Carolina using Bourne shell scripts written by Steve Bellovin. The public release of news was in the form of conventional compiled software, written by Steve Daniel and Truscott. In 1980, Usenet was connected to ARPANET through , which had connections to both Usenet and ARPANET. Mary Ann Horton, the graduate student who set up the connection, began "feeding mailing lists from the ARPANET into Usenet" with the "fa" ("From ARPANET") identifier. Usenet gained 50 member sites in its first year, including Reed College, University of Oklahoma, and Bell Labs, and the number of people using the network increased dramatically; however, it was still a while longer before Usenet users could contribute to ARPANET.
Network
UUCP networks spread quickly due to the lower costs involved, and the ability to use existing leased lines, X.25 links or even ARPANET connections. By 1983, thousands of people participated from more than 500 hosts, mostly universities and Bell Labs sites but also a growing number of Unix-related companies; the number of hosts nearly doubled to 940 in 1984. More than 100 newsgroups existed, more than 20 devoted to Unix and other computer-related topics, and at least a third to recreation. As the mesh of UUCP hosts rapidly expanded, it became desirable to distinguish the Usenet subset from the overall network. A vote was taken at the 1982 USENIX conference to choose a new name. The name Usenet was retained, but it was established that it only applied to news. The name UUCPNET became the common name for the overall network.
In addition to UUCP, early Usenet traffic was also exchanged with FidoNet and other dial-up BBS networks. By the mid-1990s there were almost 40,000 FidoNet systems in operation, and it was possible to communicate with millions of users around the world, with only local telephone service. Widespread use of Usenet by the BBS community was facilitated by the introduction of UUCP feeds made possible by MS-DOS implementations of UUCP, such as UFGATE (UUCP to FidoNet Gateway), FSUUCP and UUPC. In 1986, RFC 977 provided the Network News Transfer Protocol (NNTP) specification for distribution of Usenet articles over TCP/IP as a more flexible alternative to informal Internet transfers of UUCP traffic. Since the Internet boom of the 1990s, almost all Usenet distribution is over NNTP.
Software
Early versions of Usenet used Duke's A News software, designed for one or two articles a day. Matt Glickman and Horton at Berkeley produced an improved version called B News that could handle the rising traffic (about 50 articles a day as of late 1983). With a message format that offered compatibility with Internet mail and improved performance, it became the dominant server software. C News, developed by Geoff Collyer and Henry Spencer at the University of Toronto, was comparable to B News in features but offered considerably faster processing. In the early 1990s, InterNetNews by Rich Salz was developed to take advantage of the continuous message flow made possible by NNTP versus the batched store-and-forward design of UUCP. Since that time INN development has continued, and other news server software has also been developed.
Public venue
Usenet was the first Internet community and the place for many of the most important public developments in the pre-commercial Internet. It was the place where Tim Berners-Lee announced the launch of the World Wide Web, where Linus Torvalds announced the Linux project, and where Marc Andreessen announced the creation of the Mosaic browser and the introduction of the image tag, which revolutionized the World Wide Web by turning it into a graphical medium.
Internet jargon and history
Many jargon terms now in common use on the Internet originated or were popularized on Usenet. Likewise, many conflicts which later spread to the rest of the Internet, such as the ongoing difficulties over spamming, began on Usenet.
Decline
Sascha Segan of PC Magazine said in 2008 that "Usenet has been dying for years". Segan said that some people pointed to the Eternal September in 1993 as the beginning of Usenet's decline, when AOL began offering Usenet access. He argues that when users began putting large (non-text) files on Usenet by the late 1990s, Usenet disk space and traffic increased correspondingly. Internet service providers questioned why they needed to host binary articles.
AOL discontinued Usenet access in 2005. In May 2010, Duke University, whose implementation had started Usenet more than 30 years earlier, decommissioned its Usenet server, citing low usage and rising costs. On February 4, 2011, the Usenet news service link at the University of North Carolina at Chapel Hill (news.unc.edu) was retired after 32 years.
In response, John Biggs of TechCrunch said "As long as there are folks who think a command line is better than a mouse, the original text-only social network will live on". While there are still some active text newsgroups on Usenet, the system is now primarily used to share large files between users, and the underlying technology of Usenet remains unchanged.
Usenet traffic changes
Over time, the amount of Usenet traffic has steadily increased. the number of all text posts made in all Big-8 newsgroups averaged 1,800 new messages every hour, with an average of 25,000 messages per day. However, these averages are minuscule in comparison to the traffic in the binary groups. Much of this traffic increase reflects not an increase in discrete users or newsgroup discussions, but instead the combination of massive automated spamming and an increase in the use of newsgroups in which large files are often posted publicly. A small sampling of the change (measured in feed size per day) follows:
In 2008, Verizon Communications, Time Warner Cable and Sprint Nextel signed an agreement with Attorney General of New York Andrew Cuomo to shut down access to sources of child pornography. Time Warner Cable stopped offering access to Usenet. Verizon reduced its access to the "Big 8" hierarchies. Sprint stopped access to the alt.* hierarchies. AT&T stopped access to the hierarchies. Cuomo never specifically named Usenet in his anti-child pornography campaign. David DeJean of PC World said that some worry that the ISPs used Cuomo's campaign as an excuse to end portions of Usenet access, as it is costly for the Internet service providers and not in high demand by customers. In 2008 AOL, which no longer offered Usenet access, and the four providers that responded to the Cuomo campaign were the five largest Internet service providers in the United States; they had more than 50% of the U.S. ISP market share. On June 8, 2009, AT&T announced that it would no longer provide access to the Usenet service as of July 15, 2009.
AOL announced that it would discontinue its integrated Usenet service in early 2005, citing the growing popularity of weblogs, chat forums and on-line conferencing. The AOL community had a tremendous role in popularizing Usenet some 11 years earlier.
In August 2009, Verizon announced that it would discontinue access to Usenet on September 30, 2009. JANET announced it would discontinue Usenet service, effective July 31, 2010, citing Google Groups as an alternative. Microsoft announced that it would discontinue support for its public newsgroups (msnews.microsoft.com) from June 1, 2010, offering web forums as an alternative.
Primary reasons cited for the discontinuance of Usenet service by general ISPs include the decline in volume of actual readers due to competition from blogs, along with cost and liability concerns of increasing proportion of traffic devoted to file-sharing and spam on unused or discontinued groups.
Some ISPs did not include pressure from Cuomo's campaign against child pornography as one of their reasons for dropping Usenet feeds as part of their services. ISPs Cox and Atlantic Communications resisted the 2008 trend but both did eventually drop their respective Usenet feeds in 2010.
Archives
Public archives of Usenet articles have existed since the early days of Usenet, such as the system created by Kenneth Almquist in late 1982. Distributed archiving of Usenet posts was suggested in November 1982 by Scott Orshan, who proposed that "Every site should keep all the articles it posted, forever." Also in November of that year, Rick Adams responded to a post asking "Has anyone archived netnews, or does anyone plan to?" by stating that he was, "afraid to admit it, but I started archiving most 'useful' newsgroups as of September 18." In June 1982, Gregory G. Woodbury proposed an "automatic access to archives" system that consisted of "automatic answering of fixed-format messages to a special mail recipient on specified machines."
In 1985, two news archiving systems and one RFC were posted to the Internet. The first system, called keepnews, by Mark M. Swenson of the University of Arizona, was described as "a program that attempts to provide a sane way of extracting and keeping information that comes over Usenet." The main advantage of this system was to allow users to mark articles as worthwhile to retain. The second system, YA News Archiver by Chuq Von Rospach, was similar to keepnews, but was "designed to work with much larger archives where the wonderful quadratic search time feature of the Unix ... becomes a real problem." Von Rospach in early 1985 posted a detailed RFC for "archiving and accessing usenet articles with keyword lookup." This RFC described a program that could "generate and
maintain an archive of Usenet articles and allow looking up articles based on the article-id, subject lines, or keywords pulled out of the article itself." Also included was C code for the internal data structure of the system.
The desire to have a full text search index of archived news articles is not new either, one such request having been made in April 1991 by Alex Martelli who sought to "build some sort of keyword index for [the news archive]." In early May, Martelli posted a summary of his responses to Usenet, noting that the "most popular suggestion award must definitely go to 'lq-text' package, by Liam Quin, recently posted in alt.sources."
The Alt Sex Stories Text Repository (ASSTR) site archives and indexes erotic and pornographic stories posted to the Usenet group alt.sex.stories.
The archiving of Usenet has led to fears of loss of privacy. An archive simplifies ways to profile people. This has partly been countered with the introduction of the X-No-Archive: Yes header, which is itself controversial.
Archives by Google Groups and Deja News
Web-based archiving of Usenet posts began in March 1995 at Deja News with a very large, searchable database. In February 2001, this database was acquired by Google; Google had begun archiving Usenet posts for itself starting in the second week of August 2000.
Google Groups hosts an archive of Usenet posts dating back to May 1981. The earliest posts, which date from May 1981 to June 1991, were donated to Google by the University of Western Ontario with the help of David Wiseman and others, and were originally archived by Henry Spencer at the University of Toronto's Zoology department. The archives for late 1991 through early 1995 were provided by Kent Landfield from the NetNews CD series and Jürgen Christoffel from GMD.
Google has been criticized by Vice and Wired contributors as well as former employees for its stewardship of the archive and for breaking its search functionality.
As of January 2024, Google Groups carries a header notice, saying:
An explanatory page adds:
| Technology | Internet | null |
18589032 | https://en.wikipedia.org/wiki/Particle%20accelerator | Particle accelerator | A particle accelerator is a machine that uses electromagnetic fields to propel charged particles to very high speeds and energies to contain them in well-defined beams. Small accelerators are used for fundamental research in particle physics. Accelerators are also used as synchrotron light sources for the study of condensed matter physics. Smaller particle accelerators are used in a wide variety of applications, including particle therapy for oncological purposes, radioisotope production for medical diagnostics, ion implanters for the manufacturing of semiconductors, and accelerator mass spectrometers for measurements of rare isotopes such as radiocarbon.
Large accelerators include the Relativistic Heavy Ion Collider at Brookhaven National Laboratory in New York, and the largest accelerator, the Large Hadron Collider near Geneva, Switzerland, operated by CERN. It is a collider accelerator, which can accelerate two beams of protons to an energy of 6.5 TeV and cause them to collide head-on, creating center-of-mass energies of 13 TeV. There are more than 30,000 accelerators in operation around the world.
There are two basic classes of accelerators: electrostatic and electrodynamic (or electromagnetic) accelerators. Electrostatic particle accelerators use static electric fields to accelerate particles. The most common types are the Cockcroft–Walton generator and the Van de Graaff generator. A small-scale example of this class is the cathode-ray tube in an ordinary old television set. The achievable kinetic energy for particles in these devices is determined by the accelerating voltage, which is limited by electrical breakdown. Electrodynamic or electromagnetic accelerators, on the other hand, use changing electromagnetic fields (either magnetic induction or oscillating radio frequency fields) to accelerate particles. Since in these types the particles can pass through the same accelerating field multiple times, the output energy is not limited by the strength of the accelerating field. This class, which was first developed in the 1920s, is the basis for most modern large-scale accelerators.
Rolf Widerøe, Gustav Ising, Leó Szilárd, Max Steenbeck, and Ernest Lawrence are considered pioneers of this field, having conceived and built the first operational linear particle accelerator, the betatron, as well as the cyclotron. Because the target of the particle beams of early accelerators was usually the atoms of a piece of matter, with the goal being to create collisions with their nuclei in order to investigate nuclear structure, accelerators were commonly referred to as atom smashers in the 20th century. The term persists despite the fact that many modern accelerators create collisions between two subatomic particles, rather than a particle and an atomic nucleus.
Uses
Beams of high-energy particles are useful for fundamental and applied research in the sciences and also in many technical and industrial fields unrelated to fundamental research. There are approximately 30,000 accelerators worldwide; of these, only about 1% are research machines with energies above 1 GeV, while about 44% are for radiotherapy, 41% for ion implantation, 9% for industrial processing and research, and 4% for biomedical and other low-energy research.
Particle physics
For the most basic inquiries into the dynamics and structure of matter, space, and time, physicists seek the simplest kinds of interactions at the highest possible energies. These typically entail particle energies of many GeV, and interactions of the simplest kinds of particles: leptons (e.g. electrons and positrons) and quarks for the matter, or photons and gluons for the field quanta. Since isolated quarks are experimentally unavailable due to color confinement, the simplest available experiments involve the interactions of, first, leptons with each other, and second, of leptons with nucleons, which are composed of quarks and gluons. To study the collisions of quarks with each other, scientists resort to collisions of nucleons, which at high energy may be usefully considered as essentially 2-body interactions of the quarks and gluons of which they are composed. This elementary particle physicists tend to use machines creating beams of electrons, positrons, protons, and antiprotons, interacting with each other or with the simplest nuclei (e.g., hydrogen or deuterium) at the highest possible energies, generally hundreds of GeV or more.
The largest and highest-energy particle accelerator used for elementary particle physics is the Large Hadron Collider (LHC) at CERN, operating since 2009.
Nuclear physics and isotope production
Nuclear physicists and cosmologists may use beams of bare atomic nuclei, stripped of electrons, to investigate the structure, interactions, and properties of the nuclei themselves, and of condensed matter at extremely high temperatures and densities, such as might have occurred in the first moments of the Big Bang. These investigations often involve collisions of heavy nucleiof atoms like iron or goldat energies of several GeV per nucleon. The largest such particle accelerator is the Relativistic Heavy Ion Collider (RHIC) at Brookhaven National Laboratory.
Particle accelerators can also produce proton beams, which can produce proton-rich medical or research isotopes as opposed to the neutron-rich ones made in fission reactors; however, recent work has shown how to make 99Mo, usually made in reactors, by accelerating isotopes of hydrogen, although this method still requires a reactor to produce tritium. An example of this type of machine is LANSCE at Los Alamos National Laboratory.
Synchrotron radiation
Electrons propagating through a magnetic field emit very bright and coherent photon beams via synchrotron radiation. It has numerous uses in the study of atomic structure, chemistry, condensed matter physics, biology, and technology. A large number of synchrotron light sources exist worldwide. Examples in the U.S. are SSRL at SLAC National Accelerator Laboratory, APS at Argonne National Laboratory, ALS at Lawrence Berkeley National Laboratory, and NSLS-II at Brookhaven National Laboratory. In Europe, there are MAX IV in Lund, Sweden, BESSY in Berlin, Germany, Diamond in Oxfordshire, UK, ESRF in Grenoble, France, the latter has been used to extract detailed 3-dimensional images of insects trapped in amber.
Free-electron lasers (FELs) are a special class of light sources based on synchrotron radiation that provides shorter pulses with higher temporal coherence. A specially designed FEL is the most brilliant source of x-rays in the observable universe. The most prominent examples are the LCLS in the U.S. and European XFEL in Germany. More attention is being drawn towards soft x-ray lasers, which together with pulse shortening opens up new methods for attosecond science. Apart from x-rays, FELs are used to emit terahertz light, e.g. FELIX in Nijmegen, Netherlands, TELBE in Dresden, Germany and NovoFEL in Novosibirsk, Russia.
Thus there is a great demand for electron accelerators of moderate (GeV) energy, high intensity and high beam quality to drive light sources.
Low-energy machines and particle therapy
Everyday examples of particle accelerators are cathode ray tubes found in television sets and X-ray generators. These low-energy accelerators use a single pair of electrodes with a DC voltage of a few thousand volts between them. In an X-ray generator, the target itself is one of the electrodes. A low-energy particle accelerator called an ion implanter is used in the manufacture of integrated circuits.
At lower energies, beams of accelerated nuclei are also used in medicine as particle therapy, for the treatment of cancer.
DC accelerator types capable of accelerating particles to speeds sufficient to cause nuclear reactions are Cockcroft–Walton generators or voltage multipliers, which convert AC to high voltage DC, or Van de Graaff generators that use static electricity carried by belts.
Radiation sterilization of medical devices
Electron beam processing is commonly used for sterilization. Electron beams are an on-off technology that provide a much higher dose rate than gamma or X-rays emitted by radioisotopes like cobalt-60 (60Co) or caesium-137 (137Cs). Due to the higher dose rate, less exposure time is required and polymer degradation is reduced. Because electrons carry a charge, electron beams are less penetrating than both gamma and X-rays.
Electrostatic particle accelerators
Historically, the first accelerators used simple technology of a single static high voltage to accelerate charged particles. The charged particle was accelerated through an evacuated tube with an electrode at either end, with the static potential across it. Since the particle passed only once through the potential difference, the output energy was limited to the accelerating voltage of the machine. While this method is still extremely popular today, with the electrostatic accelerators greatly out-numbering any other type, they are more suited to lower energy studies owing to the practical voltage limit of about 1 MV for air insulated machines, or 30 MV when the accelerator is operated in a tank of pressurized gas with high dielectric strength, such as sulfur hexafluoride. In a tandem accelerator the potential is used twice to accelerate the particles, by reversing the charge of the particles while they are inside the terminal. This is possible with the acceleration of atomic nuclei by using anions (negatively charged ions), and then passing the beam through a thin foil to strip electrons off the anions inside the high voltage terminal, converting them to cations (positively charged ions), which are accelerated again as they leave the terminal.
The two main types of electrostatic accelerator are the Cockcroft–Walton accelerator, which uses a diode-capacitor voltage multiplier to produce high voltage, and the Van de Graaff accelerator, which uses a moving fabric belt to carry charge to the high voltage electrode. Although electrostatic accelerators accelerate particles along a straight line, the term linear accelerator is more often used for accelerators that employ oscillating rather than static electric fields.
Electrodynamic (electromagnetic) particle accelerators
Due to the high voltage ceiling imposed by electrical discharge, in order to accelerate particles to higher energies, techniques involving dynamic fields rather than static fields are used. Electrodynamic acceleration can arise from either of two mechanisms: non-resonant magnetic induction, or resonant circuits or cavities excited by oscillating radio frequency (RF) fields.
Electrodynamic accelerators can be linear, with particles accelerating in a straight line, or circular, using magnetic fields to bend particles in a roughly circular orbit.
Magnetic induction accelerators
Magnetic induction accelerators accelerate particles by induction from an increasing magnetic field, as if the particles were the secondary winding in a transformer. The increasing magnetic field creates a circulating electric field which can be configured to accelerate the particles. Induction accelerators can be either linear or circular.
Linear induction accelerators
Linear induction accelerators utilize ferrite-loaded, non-resonant induction cavities. Each cavity can be thought of as two large washer-shaped disks connected by an outer cylindrical tube. Between the disks is a ferrite toroid. A voltage pulse applied between the two disks causes an increasing magnetic field which inductively couples power into the charged particle beam.
The linear induction accelerator was invented by Christofilos in the 1960s. Linear induction accelerators are capable of accelerating very high beam currents (>1000 A) in a single short pulse. They have been used to generate X-rays for flash radiography (e.g. DARHT at LANL), and have been considered as particle injectors for magnetic confinement fusion and as drivers for free electron lasers.
Betatrons
The Betatron is a circular magnetic induction accelerator, invented by Donald Kerst in 1940 for accelerating electrons. The concept originates ultimately from Norwegian-German scientist Rolf Widerøe. These machines, like synchrotrons, use a donut-shaped ring magnet (see below) with a cyclically increasing B field, but accelerate the particles by induction from the increasing magnetic field, as if they were the secondary winding in a transformer, due to the changing magnetic flux through the orbit.
Achieving constant orbital radius while supplying the proper accelerating electric field requires that the magnetic flux linking the orbit be somewhat independent of the magnetic field on the orbit, bending the particles into a constant radius curve. These machines have in practice been limited by the large radiative losses suffered by the electrons moving at nearly the speed of light in a relatively small radius orbit.
Linear accelerators
In a linear particle accelerator (linac), particles are accelerated in a straight line with a target of interest at one end. They are often used to provide an initial low-energy kick to particles before they are injected into circular accelerators. The longest linac in the world is the Stanford Linear Accelerator, SLAC, which is long. SLAC was originally an electron–positron collider but is now a X-ray Free-electron laser.
Linear high-energy accelerators use a linear array of plates (or drift tubes) to which an alternating high-energy field is applied. As the particles approach a plate they are accelerated towards it by an opposite polarity charge applied to the plate. As they pass through a hole in the plate, the polarity is switched so that the plate now repels them and they are now accelerated by it towards the next plate. Normally a stream of "bunches" of particles are accelerated, so a carefully controlled AC voltage is applied to each plate to continuously repeat this process for each bunch.
As the particles approach the speed of light the switching rate of the electric fields becomes so high that they operate at radio frequencies, and so microwave cavities are used in higher energy machines instead of simple plates.
Linear accelerators are also widely used in medicine, for radiotherapy and radiosurgery. Medical grade linacs accelerate electrons using a klystron and a complex bending magnet arrangement which produces a beam of energy . The electrons can be used directly or they can be collided with a target to produce a beam of X-rays. The reliability, flexibility and accuracy of the radiation beam produced has largely supplanted the older use of cobalt-60 therapy as a treatment tool.
Circular or cyclic RF accelerators
In the circular accelerator, particles move in a circle until they reach enough energy. The particle track is typically bent into a circle using electromagnets. The advantage of circular accelerators over linear accelerators (linacs) is that the ring topology allows continuous acceleration, as the particle can transit indefinitely. Another advantage is that a circular accelerator is smaller than a linear accelerator of comparable power (i.e. a linac would have to be extremely long to have the equivalent power of a circular accelerator).
Depending on the energy and the particle being accelerated, circular accelerators suffer a disadvantage in that the particles emit synchrotron radiation. When any charged particle is accelerated, it emits electromagnetic radiation and secondary emissions. As a particle traveling in a circle is always accelerating towards the center of the circle, it continuously radiates towards the tangent of the circle. This radiation is called synchrotron light and depends highly on the mass of the accelerating particle. For this reason, many high energy electron accelerators are linacs. Certain accelerators (synchrotrons) are however built specially for producing synchrotron light (X-rays).
Since the special theory of relativity requires that matter always travels slower than the speed of light in vacuum, in high-energy accelerators, as the energy increases the particle speed approaches the speed of light as a limit, but never attains it. Therefore, particle physicists do not generally think in terms of speed, but rather in terms of a particle's energy or momentum, usually measured in electron volts (eV). An important principle for circular accelerators, and particle beams in general, is that the curvature of the particle trajectory is proportional to the particle charge and to the magnetic field, but inversely proportional to the (typically relativistic) momentum.
Cyclotrons
The earliest operational circular accelerators were cyclotrons, invented in 1929 by Ernest Lawrence at the University of California, Berkeley. Cyclotrons have a single pair of hollow D-shaped plates to accelerate the particles and a single large dipole magnet to bend their path into a circular orbit. It is a characteristic property of charged particles in a uniform and constant magnetic field B that they orbit with a constant period, at a frequency called the cyclotron frequency, so long as their speed is small compared to the speed of light c. This means that the accelerating D's of a cyclotron can be driven at a constant frequency by a RF accelerating power source, as the beam spirals outwards continuously. The particles are injected in the center of the magnet and are extracted at the outer edge at their maximum energy.
Cyclotrons reach an energy limit because of relativistic effects whereby the particles effectively become more massive, so that their cyclotron frequency drops out of sync with the accelerating RF. Therefore, simple cyclotrons can accelerate protons only to an energy of around 15 million electron volts (15 MeV, corresponding to a speed of roughly 10% of c), because the protons get out of phase with the driving electric field. If accelerated further, the beam would continue to spiral outward to a larger radius but the particles would no longer gain enough speed to complete the larger circle in step with the accelerating RF. To accommodate relativistic effects the magnetic field needs to be increased to higher radii as is done in isochronous cyclotrons. An example of an isochronous cyclotron is the PSI Ring cyclotron in Switzerland, which provides protons at the energy of 590 MeV which corresponds to roughly 80% of the speed of light. The advantage of such a cyclotron is the maximum achievable extracted proton current which is currently 2.2 mA. The energy and current correspond to 1.3 MW beam power which is the highest of any accelerator currently existing.
Synchrocyclotrons and isochronous cyclotrons
A classic cyclotron can be modified to increase its energy limit. The historically first approach was the synchrocyclotron, which accelerates the particles in bunches. It uses a constant magnetic field , but reduces the accelerating field's frequency so as to keep the particles in step as they spiral outward, matching their mass-dependent cyclotron resonance frequency. This approach suffers from low average beam intensity due to the bunching, and again from the need for a huge magnet of large radius and constant field over the larger orbit demanded by high energy.
The second approach to the problem of accelerating relativistic particles is the isochronous cyclotron. In such a structure, the accelerating field's frequency (and the cyclotron resonance frequency) is kept constant for all energies by shaping the magnet poles so to increase magnetic field with radius. Thus, all particles get accelerated in isochronous time intervals. Higher energy particles travel a shorter distance in each orbit than they would in a classical cyclotron, thus remaining in phase with the accelerating field. The advantage of the isochronous cyclotron is that it can deliver continuous beams of higher average intensity, which is useful for some applications. The main disadvantages are the size and cost of the large magnet needed, and the difficulty in achieving the high magnetic field values required at the outer edge of the structure.
Synchrocyclotrons have not been built since the isochronous cyclotron was developed.
Synchrotrons
To reach still higher energies, with relativistic mass approaching or exceeding the rest mass of the particles (for protons, billions of electron volts or GeV), it is necessary to use a synchrotron. This is an accelerator in which the particles are accelerated in a ring of constant radius. An immediate advantage over cyclotrons is that the magnetic field need only be present over the actual region of the particle orbits, which is much narrower than that of the ring. (The largest cyclotron built in the US had a magnet pole, whereas the diameter of synchrotrons such as the LEP and LHC is nearly 10 km. The aperture of the two beams of the LHC is of the order of a centimeter.) The LHC contains 16 RF cavities, 1232 superconducting dipole magnets for beam steering, and 24 quadrupoles for beam focusing. Even at this size, the LHC is limited by its ability to steer the particles without them going adrift. This limit is theorized to occur at 14 TeV.
However, since the particle momentum increases during acceleration, it is necessary to turn up the magnetic field B in proportion to maintain constant curvature of the orbit. In consequence, synchrotrons cannot accelerate particles continuously, as cyclotrons can, but must operate cyclically, supplying particles in bunches, which are delivered to a target or an external beam in beam "spills" typically every few seconds.
Since high energy synchrotrons do most of their work on particles that are already traveling at nearly the speed of light c, the time to complete one orbit of the ring is nearly constant, as is the frequency of the RF cavity resonators used to drive the acceleration.
In modern synchrotrons, the beam aperture is small and the magnetic field does not cover the entire area of the particle orbit as it does for a cyclotron, so several necessary functions can be separated. Instead of one huge magnet, one has a line of hundreds of bending magnets, enclosing (or enclosed by) vacuum connecting pipes. The design of synchrotrons was revolutionized in the early 1950s with the discovery of the strong focusing concept. The focusing of the beam is handled independently by specialized quadrupole magnets, while the acceleration itself is accomplished in separate RF sections, rather similar to short linear accelerators. Also, there is no necessity that cyclic machines be circular, but rather the beam pipe may have straight sections between magnets where beams may collide, be cooled, etc. This has developed into an entire separate subject, called "beam physics" or "beam optics".
More complex modern synchrotrons such as the Tevatron, LEP, and LHC may deliver the particle bunches into storage rings of magnets with a constant magnetic field, where they can continue to orbit for long periods for experimentation or further acceleration. The highest-energy machines such as the Tevatron and LHC are actually accelerator complexes, with a cascade of specialized elements in series, including linear accelerators for initial beam creation, one or more low energy synchrotrons to reach intermediate energy, storage rings where beams can be accumulated or "cooled" (reducing the magnet aperture required and permitting tighter focusing; see beam cooling), and a last large ring for final acceleration and experimentation.
Electron synchrotrons
Circular electron accelerators fell somewhat out of favor for particle physics around the time that SLAC's linear particle accelerator was constructed, because their synchrotron losses were considered economically prohibitive and because their beam intensity was lower than for the unpulsed linear machines. The Cornell Electron Synchrotron, built at low cost in the late 1970s, was the first in a series of high-energy circular electron accelerators built for fundamental particle physics, the last being LEP, built at CERN, which was used from 1989 until 2000.
A large number of electron synchrotrons have been built in the past two decades, as part of synchrotron light sources that emit ultraviolet light and X rays; see below.
Synchrotron radiation sources
Some circular accelerators have been built to deliberately generate radiation (called synchrotron light) as X-rays also called synchrotron radiation, for example the Diamond Light Source which has been built at the Rutherford Appleton Laboratory in England or the Advanced Photon Source at Argonne National Laboratory in Illinois, USA. High-energy X-rays are useful for X-ray spectroscopy of proteins or X-ray absorption fine structure (XAFS), for example.
Synchrotron radiation is more powerfully emitted by lighter particles, so these accelerators are invariably electron accelerators. Synchrotron radiation allows for better imaging as researched and developed at SLAC's SPEAR.
Fixed-field alternating gradient accelerators
Fixed-Field Alternating Gradient accelerators (FFA)s, in which a magnetic field which is fixed in time, but with a radial variation to achieve strong focusing, allows the beam to be accelerated with a high repetition rate but in a much smaller radial spread than in the cyclotron case. Isochronous FFAs, like isochronous cyclotrons, achieve continuous beam operation, but without the need for a huge dipole bending magnet covering the entire radius of the orbits. Some new developments in FFAs are covered in.
Rhodotron
A Rhodotron is an industrial electron accelerator first proposed in 1987 by J. Pottier of the French Atomic
Energy Agency (CEA), manufactured by Belgian company Ion Beam Applications. It accelerates electrons by recirculating them across the diameter of a cylinder-shaped radiofrequency cavity. A Rhodotron has an electron gun, which emits an electron beam that is attracted to a pillar in the center of the cavity. The pillar has holes the electrons can pass through. The electron beam passes through the pillar via one of these holes and then travels through a hole in the wall of the cavity, and meets a bending magnet, the beam is then bent and sent back into the cavity, to another hole in the pillar, the electrons then again go across the pillar and pass though another part of the wall of the cavity and into another bending magnet, and so on, gradually increasing the energy of the beam until it is allowed to exit the cavity for use. The cylinder and pillar may be lined with copper on the inside.
History
Ernest Lawrence's first cyclotron was a mere 4 inches (100 mm) in diameter. Later, in 1939, he built a machine with a 60-inch diameter pole face, and planned one with a 184-inch diameter in 1942, which was, however, taken over for World War II-related work connected with uranium isotope separation; after the war it continued in service for research and medicine over many years.
The first large proton synchrotron was the Cosmotron at Brookhaven National Laboratory, which accelerated protons to about 3 GeV (1953–1968). The Bevatron at Berkeley, completed in 1954, was specifically designed to accelerate protons to enough energy to create antiprotons, and verify the particle–antiparticle symmetry of nature, then only theorized. The Alternating Gradient Synchrotron (AGS) at Brookhaven (1960–) was the first large synchrotron with alternating gradient, "strong focusing" magnets, which greatly reduced the required aperture of the beam, and correspondingly the size and cost of the bending magnets. The Proton Synchrotron, built at CERN (1959–), was the first major European particle accelerator and generally similar to the AGS.
The Stanford Linear Accelerator, SLAC, became operational in 1966, accelerating electrons to 30 GeV in a 3 km long waveguide, buried in a tunnel and powered by hundreds of large klystrons. It is still the largest linear accelerator in existence, and has been upgraded with the addition of storage rings and an electron-positron collider facility. It is also an X-ray and UV synchrotron photon source.
The Fermilab Tevatron has a ring with a beam path of . It has received several upgrades, and has functioned as a proton-antiproton collider until it was shut down due to budget cuts on September 30, 2011. The largest circular accelerator ever built was the LEP synchrotron at CERN with a circumference 26.6 kilometers, which was an electron/positron collider. It achieved an energy of 209 GeV before it was dismantled in 2000 so that the tunnel could be used for the Large Hadron Collider (LHC). The LHC is a proton collider, and currently the world's largest and highest-energy accelerator, achieving 6.5 TeV energy per beam (13 TeV in total).
The aborted Superconducting Super Collider (SSC) in Texas would have had a circumference of 87 km. Construction was started in 1991, but abandoned in 1993. Very large circular accelerators are invariably built in tunnels a few metres wide to minimize the disruption and cost of building such a structure on the surface, and to provide shielding against intense secondary radiations that occur, which are extremely penetrating at high energies.
Current accelerators such as the Spallation Neutron Source, incorporate superconducting cryomodules. The Relativistic Heavy Ion Collider, and Large Hadron Collider also make use of superconducting magnets and RF cavity resonators to accelerate particles.
Targets
The output of a particle accelerator can generally be directed towards multiple lines of experiments, one at a given time, by means of a deviating electromagnet. This makes it possible to operate multiple experiments without needing to move things around or shutting down the entire accelerator beam. Except for synchrotron radiation sources, the purpose of an accelerator is to generate high-energy particles for interaction with matter.
This is usually a fixed target, such as the phosphor coating on the back of the screen in the case of a television tube; a piece of uranium in an accelerator designed as a neutron source; or a tungsten target for an X-ray generator. In a linac, the target is simply fitted to the end of the accelerator. The particle track in a cyclotron is a spiral outwards from the centre of the circular machine, so the accelerated particles emerge from a fixed point as for a linear accelerator.
For synchrotrons, the situation is more complex. Particles are accelerated to the desired energy. Then, a fast acting dipole magnet is used to switch the particles out of the circular synchrotron tube and towards the target.
A variation commonly used for particle physics research is a collider, also called a storage ring collider. Two circular synchrotrons are built in close proximityusually on top of each other and using the same magnets (which are then of more complicated design to accommodate both beam tubes). Bunches of particles travel in opposite directions around the two accelerators and collide at intersections between them. This can increase the energy enormously; whereas in a fixed-target experiment the energy available to produce new particles is proportional to the square root of the beam energy, in a collider the available energy is linear.
Detectors NA
The detectors gather clues about the particles including their speed and charge. Using these, the scientists can actually work on the particle. The process of detection is very complex it requires strong electromagnets and accelerators to generate enough usable information.
Higher energies
At present the highest energy accelerators are all circular colliders, but both hadron accelerators and electron accelerators are running into limits. Higher energy hadron and ion cyclic accelerators will require accelerator tunnels of larger physical size due to the increased beam rigidity.
For cyclic electron accelerators, a limit on practical bend radius is placed by synchrotron radiation losses and the next generation will probably be linear accelerators 10 times the current length. An example of such a next generation electron accelerator is the proposed 40 km long International Linear Collider.
It is believed that plasma wakefield acceleration in the form of electron-beam "afterburners" and standalone laser pulsers might be able to provide dramatic increases in efficiency over RF accelerators within two to three decades. In plasma wakefield accelerators, the beam cavity is filled with a plasma (rather than vacuum). A short pulse of electrons or laser light either constitutes or immediately precedes the particles that are being accelerated. The pulse disrupts the plasma, causing the charged particles in the plasma to integrate into and move toward the rear of the bunch of particles that are being accelerated. This process transfers energy to the particle bunch, accelerating it further, and continues as long as the pulse is coherent.
Energy gradients as steep as 200 GeV/m have been achieved over millimeter-scale distances using laser pulsers and gradients approaching 1 GeV/m are being produced on the multi-centimeter-scale with electron-beam systems, in contrast to a limit of about 0.1 GeV/m for radio-frequency acceleration alone. Existing electron accelerators such as SLAC could use electron-beam afterburners to greatly increase the energy of their particle beams, at the cost of beam intensity. Electron systems in general can provide tightly collimated, reliable beams; laser systems may offer more power and compactness. Thus, plasma wakefield accelerators could be used – if technical issues can be resolved – to both increase the maximum energy of the largest accelerators and to bring high energies into university laboratories and medical centres.
Higher than 0.25 GeV/m gradients have been achieved by a dielectric laser accelerator, which may present another viable approach to building compact high-energy accelerators. Using femtosecond duration laser pulses, an electron accelerating gradient 0.69 GeV/m was recorded for dielectric laser accelerators. Higher gradients of the order of are anticipated after further optimizations.
Advanced Accelerator Concepts
Advanced Accelerator Concepts encompasses methods of beam acceleration with gradients beyond state of the art in operational facilities. This includes diagnostics methods, timing technology, special needs for injectors, beam matching, beam dynamics and development of adequate simulations. Workshops dedicated to this subject are being held in the US (alternating locations) and in Europe, mostly on Isola d'Elba.
The series of Advanced Accelerator Concepts Workshops, held in the US, started as an international series in 1982.
The European Advanced Accelerator Concepts Workshop series started in 2019.
Topics related to Advanced Accelerator Concepts:
Laser Plasma Acceleration of electrons and positrons
Laser and High-Gradient Structure-Based Acceleration
Beam-Driven Acceleration
Laser-Plasma Acceleration of Ions
Beam Sources such as electron gun, Monitoring, and Control. See Accelerator physics
Computer simulation for Accelerator Physics
Laser technology for particle acceleration
Electromagnetic radiation Generation
Muon collider
According to the Inverse scattering problem, any mechanism by which a particle produces radiation (where kinetic energy of the particle is transferred to the electromagnetic field), can be inverted such that the same radiation mechanism leads to the acceleration of the particle (energy of the radiation field is transferred to kinetic energy of the particle). The opposite is also true, any acceleration mechanism can be inverted to deposit the energy of the particle into a decelerating field, like in a kinetic energy recovery system. This is the idea enabling an energy recovery linac.
This principle, which is also behind the plasma or dielectric wakefield accelerrators, led to a few other interesting developments in advanced accelerator concepts:
Cherenkov radiation led to inverse Cherenkov radiation accelerator.
Free-electron laser led to the Inverse Free-electron laser accelerator.
A laser can also be inverted to produce acceleration of electrons.
Black hole production and public safety concerns
In the future, the possibility of a black hole production at the highest energy accelerators may arise if certain predictions of superstring theory are accurate. This and other possibilities have led to public safety concerns that have been widely reported in connection with the LHC, which began operation in 2008. The various possible dangerous scenarios have been assessed as presenting "no conceivable danger" in the latest risk assessment produced by the LHC Safety Assessment Group. If black holes are produced, it is theoretically predicted that such small black holes should evaporate extremely quickly via Bekenstein–Hawking radiation, but which is as yet experimentally unconfirmed. If colliders can produce black holes, cosmic rays (and particularly ultra-high-energy cosmic rays, UHECRs) must have been producing them for eons, but they have yet to harm anybody. It has been argued that to conserve energy and momentum, any black holes created in a collision between an UHECR and local matter would necessarily be produced moving at relativistic speed with respect to the Earth, and should escape into space, as their accretion and growth rate should be very slow, while black holes produced in colliders (with components of equal mass) would have some chance of having a velocity less than Earth escape velocity, 11.2 km per sec, and would be liable to capture and subsequent growth. Yet even on such scenarios the collisions of UHECRs with white dwarfs and neutron stars would lead to their rapid destruction, but these bodies are observed to be common astronomical objects. Thus if stable micro black holes should be produced, they must grow far too slowly to cause any noticeable macroscopic effects within the natural lifetime of the solar system.
Accelerator operator
The use of advanced technologies such as superconductivity, cryogenics, and high powered radiofrequency amplifiers, as well as the presence of ionizing radiation, pose challenges for the safe operation of accelerator facilities. An accelerator operator controls the operation of a particle accelerator, adjusts operating parameters such as aspect ratio, current intensity, and position on target. They communicate with and assist accelerator maintenance personnel to ensure readiness of support systems, such as vacuum, magnets, magnetic and radiofrequency power supplies and controls, and cooling systems. Additionally, the accelerator operator maintains a record of accelerator related events.
| Physical sciences | Particle physics: General | null |
18589212 | https://en.wikipedia.org/wiki/Reverse%20osmosis | Reverse osmosis | Reverse osmosis (RO) is a water purification process that uses a semi-permeable membrane to separate water molecules from other substances. RO applies pressure to overcome osmotic pressure that favors even distributions. RO can remove dissolved or suspended chemical species as well as biological substances (principally bacteria), and is used in industrial processes and the production of potable water. RO retains the solute on the pressurized side of the membrane and the purified solvent passes to the other side. The relative sizes of the various molecules determines what passes through. "Selective" membranes reject large molecules, while accepting smaller molecules (such as solvent molecules, e.g., water).
RO is most commonly known for its use in drinking water purification from seawater, removing the salt and other effluent materials from the water molecules.
As of 2013 the world's largest RO desalination plant was in Sorek, Israel, outputting .
History
A process of osmosis through semi-permeable membranes was first observed in 1748 by Jean-Antoine Nollet. For the following 200 years, osmosis was only a laboratory phenomenon. In 1950, the University of California at Los Angeles (UCLA) first investigated osmotic desalination. Researchers at both UCLA and University of Florida desalinated seawater in the mid-1950s, but the flux was too low to be commercially viable. Sidney Loeb at UCLA and Srinivasa Sourirajan at the National Research Council of Canada, Ottawa, found techniques for making asymmetric membranes characterized by an effectively thin "skin" layer supported atop a highly porous and much thicker substrate region. John Cadotte, of Filmtec corporation, discovered that membranes with particularly high flux and low salt passage could be made by interfacial polymerization of m-phenylene diamine and trimesoyl chloride. Cadotte's patent on this process was the subject of litigation and expired. Almost all commercial RO membrane is now made by this method. By 2019, approximately 16,000 desalination plants operated around the world, producing around . Around half of this capacity was in the Middle East and North Africa region.
In 1977 Cape Coral, Florida became the first US municipality to use RO at scale, with an initial operating capacity of 11.35 million liters (3 million US gal) per day. By 1985, rapid growth led the city to operate the world's largest low-pressure RO plant, producing 56.8 million liters (15 million US gal) per day (MGD).
Osmosis
In (forward) osmosis, the solvent moves from an area of low solute concentration (high water potential), through a membrane, to an area of high solute concentration (low water potential). The driving force for the movement of the solvent is the reduction in the Gibbs free energy of the system in which the difference in solvent concentration between the sides of a membrane is reduced. This is called osmotic pressure. It reduces as the solvent moves into the more concentrated solution. Applying an external pressure to reverse the natural flow of pure solvent, thus, is reverse osmosis. The process is similar to other membrane technology applications.
RO differs from filtration in that the mechanism of fluid flow is reversed, as the solvent crosses membrane, leaving the solute behind. The predominant removal mechanism in membrane filtration is straining, or size exclusion, where the pores are 0.01 micrometers or larger, so the process can theoretically achieve perfect efficiency regardless of parameters such as the solution's pressure and concentration. RO instead involves solvent diffusion across a membrane that is either nonporous or uses nanofiltration with pores 0.001 micrometers in size. The predominant removal mechanism is from differences in solubility or diffusivity, and the process is dependent on pressure, solute concentration, and other conditions.
RO requires pressure between 2–17 bar (30–250 psi) for fresh and brackish water, and 40–82 bar (600–1200 psi) for seawater. Seawater has around 27 bar (390 psi) natural osmotic pressure that must be overcome. As for their energy consumption, seawater RO systems typically require 2.9-5.5 kWh/m3 , although state-of-the-art systems are around 2.3 kWh/m3 .
Membrane pore sizes vary from 0.1 to 5,000 nm. Particle filtration removes particles of 1 μm or larger. Microfiltration removes particles of 50 nm or larger. Ultrafiltration removes particles of roughly 3 nm or larger. Nanofiltration removes particles of 1 nm or larger. RO is in the final category of membrane filtration, hyperfiltration, and removes particles larger than ~0.2 nm.
Fresh water applications
Drinking water purification
Around the world, household drinking water purification systems, including an RO step, are commonly used for improving water for drinking and cooking.
Such systems typically include these steps:
a sediment filter to trap particles, including rust and calcium carbonate
a second sediment filter with smaller pores
an activated carbon filter to trap organic chemicals and chlorine, which degrades certain types of thin-film composite membrane
an RO thin-film composite membrane
an ultraviolet lamp for sterilizing any microbes that survive RO
a second carbon filter to capture chemicals that survive RO
In some systems, the carbon prefilter is replaced by a cellulose triacetate (CTA) membrane. CTA is a paper by-product membrane bonded to a synthetic layer that allows contact with chlorine in the water. These require a small amount of chlorine in the water source to prevent bacteria from forming on it. The typical rejection rate for CTA membranes is 85–95%.
The cellulose triacetate membrane rots unless protected by chlorinated water, while the thin-film composite membrane breaks down in the presence of chlorine. The thin-film composite (TFC) membrane is made of synthetic material, and requires the chlorine to be removed before the water enters the membrane. To protect the TFC membrane elements from chlorine damage, carbon filters are used as pre-treatment. TFC membranes have a higher rejection rate of 95–98% and a longer life than CTA membranes.
To work effectively, the water feeding to these units should be under pressure (typically 280 kPa (40 psi) or greater).
Though Portable RO Water Purifiers are commercially available and extensively used in areas lacking cleaning potable water, in Europe such processing of natural mineral water (as defined by a European directive) is not allowed. In practice, a fraction of the living bacteria pass through RO through membrane imperfections or bypass the membrane entirely through leaks in seals.
Solar-powered RO
A solar-powered desalination unit produces potable water from saline water by using a photovoltaic system to supply the energy. Solar power works well for water purification in settings lacking grid electricity and can reduce operating costs and greenhouse emissions. For example, a solar-powered desalination unit designed passed tests in Australia's Northern Territory.
Sunlight's intermittent nature makes output prediction difficult without an energy storage capability. However batteries or thermal energy storage systems can provide power when the sun does not.
Military
Larger scale reverse osmosis water purification units (ROWPU) exist for military use. These have been adopted by the United States armed forces and the Canadian Forces. Some models are containerized, some are trailers, and some are themselves vehicles.
The water is treated with a polymer to initiate coagulation. Next, it is run through a multi-media filter where it undergoes primary treatment, removing turbidity. It is then pumped through a cartridge filter which is usually spiral-wound cotton. This process strips any particles larger than 5 μm and eliminates almost all turbidity.
The clarified water is then fed through a high-pressure piston pump into a series of RO vessels. 90.00–99.98% of the raw water's total dissolved solids are removed and military standards require that the result have no more than 1000–1500 parts per million by measure of electrical conductivity. It is then disinfected with chlorine.
Water and wastewater purification
RO-purified rainwater collected from storm drains is used for landscape irrigation and industrial cooling in Los Angeles and other cities.
In industry, RO removes minerals from boiler water at power plants. The water is distilled multiple times to ensure that it does not leave deposits on the machinery or cause corrosion.
RO is used to clean effluent and brackish groundwater. The effluent in larger volumes (more than 500 m3/day) is treated in a water treatment plant first, and then the effluent runs through RO. This hybrid process reduces treatment cost significantly and lengthens membrane life.
RO can be used for the production of deionized water.
In 2002, Singapore announced that a process named NEWater would be a significant part of its water plans. RO would be used to treat wastewater before discharging the effluent into reservoirs.
Food industry
Reverse osmosis is a more economical way to concentrate liquids (such as fruit juices) than conventional heat-treatment. Concentration of orange and tomato juice has advantages including a lower operating cost and the ability to avoid heat-treatment, which makes it suitable for heat-sensitive substances such as protein and enzymes.
RO is used in the dairy industry to produce whey protein powders and concentrate milk. The whey (liquid remaining after cheese manufacture) is concentrated with RO from 6% solids to 10–20% solids before ultrafiltration processing. The retentate can then be used to make whey powders, including whey protein isolate. Additionally, the permeate, which contains lactose, is concentrated by RO from 5% solids to 18–total solids to reduce crystallization and drying costs.
Although RO was once avoided in the wine industry, it is now widespread. An estimated 60 RO machines were in use in Bordeaux, France, in 2002. Known users include many of elite firms, such as Château Léoville-Las Cases.
Maple syrup production
In 1946, some maple syrup producers started using RO to remove water from sap before boiling the sap to syrup. RO allows about 75–90% of the water to be removed, reducing energy consumption and exposure of the syrup to high temperatures.
Low-alcohol beer
When beer at typical concentration is subjected to reverse osmosis, both water and alcohol pass across the membrane more readily than other components, leaving a "beer concentrate". The concentrate is then diluted with fresh water to restore the non-volatile components to their original intensity.
Hydrogen production
For small-scale hydrogen production, RO is sometimes used to prevent formation of mineral deposits on the surface of electrodes.
Aquariums
Many reef aquarium keepers use RO systems to make fish-friendly seawater. Ordinary tap water can contain excessive chlorine, chloramines, copper, nitrates, nitrites, phosphates, silicates, or other chemicals detrimental to marine organisms. Contaminants such as nitrogen and phosphates can lead to unwanted algae growth. An effective combination of both RO and deionization is popular among reef aquarium keepers, and is preferred above other water purification processes due to the low cost of ownership and operating costs. Where chlorine and chloramines are found in the water, carbon filtration is needed before RO, as common residential membranes do not address these compounds.
Freshwater aquarists also use RO to duplicate the soft waters found in many tropical waters. While many tropical fish can survive in treated tap water, breeding can be impossible. Many aquatic shops sell containers of RO water for this purpose.
Window cleaning
An increasingly popular method of cleaning windows is the "water-fed pole" system. Instead of washing windows with conventional detergent, they are scrubbed with purified water, typically containing less than 10 ppm dissolved solids, using a brush on the end of a pole wielded from ground level. RO is commonly used to purify the water.
Landfill leachate purification
Treatment with RO is limited, resulting in low recoveries on high concentration (measured with electrical conductivity) and membrane fouling. RO applicability is limited by conductivity, organics, and scaling inorganic elements such as CaSO4, Si, Fe and Ba. Low organic scaling can use two different technologies: spiral wound membrane, and (for high organic scaling, high conductivity and higher pressure (up to 90 bars)), disc tube modules with RO membranes can be used. Disc tube modules were redesigned for landfill leachate purification that is usually contaminated with organic material. Due to the cross-flow, it is given a flow booster pump that recirculates the flow over the membrane between 1.5 and 3 times before it is released as a concentrate. High velocity protects against membrane scaling and allows membrane cleaning.
Power consumption for a disc tube module system
Desalination
Areas that have limited surface water or groundwater may choose to desalinate. RO is an increasingly common method, because of its relatively low energy consumption.
Energy consumption is around , with the development of more efficient energy recovery devices and improved membrane materials. According to the International Desalination Association, for 2011, RO was used in 66% of installed desalination capacity (0.0445 of 0.0674 km3/day), and nearly all new plants. Other plants use thermal distillation methods: multiple-effect distillation, and multi-stage flash.
Sea-water RO (SWRO) desalination requires around 3 kWh/m3, much higher than those required for other forms of water supply, including RO treatment of wastewater, at 0.1 to 1 kWh/m3. Up to 50% of the seawater input can be recovered as fresh water, though lower recovery rates may reduce membrane fouling and energy consumption.
Brackish water reverse osmosis (BWRO) is the desalination of water with less salt than seawater, usually from river estuaries or saline wells. The process is substantially the same as SWRO, but requires lower pressures and less energy. Up to 80% of the feed water input can be recovered as fresh water, depending on feed salinity.
The Ashkelon desalination plant in Israel is the world's largest.
The typical single-pass SWRO system consists of:
Intake
Pretreatment
High-pressure pump (if not combined with energy recovery)
Membrane assembly
Energy recovery (if used)
Remineralisation and pH adjustment
Disinfection
Alarm/control panel
Pretreatment
Pretreatment is important when working nanofiltration membranes due to their spiral-wound design. The material is engineered to allow one-way flow. The design does not allow for backpulsing with water or air agitation to scour its surface and remove accumulated solids. Since material cannot be removed from the membrane surface, it is susceptible to fouling (loss of production capacity). Therefore, pretreatment is a necessity for any RO or nanofiltration system. Pretreatment has four major components:
Screening solids: Solids must be removed and the water treated to prevent membrane fouling by particle or biological growth, and reduce the risk of damage to high-pressure components.
Cartridge filtration: String-wound polypropylene filters are typically used to remove particles of 1–5 μm diameter.
Dosing: Oxidizing biocides, such as chlorine, are added to kill bacteria, followed by bisulfite dosing to deactivate the chlorine that can destroy a thin-film composite membrane. Biofouling inhibitors do not kill bacteria, while preventing them from growing slime on the membrane surface and plant walls.
Prefiltration pH adjustment: If the pH, hardness and the alkalinity in the feedwater result in scaling while concentrated in the reject stream, acid is dosed to maintain carbonates in their soluble carbonic acid form.
CO32− + H3O+ = HCO3− + H2O
HCO3− + H3O+ = H2CO3 + H2O
Carbonic acid cannot combine with calcium to form calcium carbonate scale. Calcium carbonate scaling tendency is estimated using the Langelier saturation index. Adding too much sulfuric acid to control carbonate scales may result in calcium sulfate, barium sulfate, or strontium sulfate scale formation on the membrane.
Prefiltration antiscalants: Scale inhibitors (also known as antiscalants) prevent formation of more scales than acid, which can only prevent formation of calcium carbonate and calcium phosphate scales. In addition to inhibiting carbonate and phosphate scales, antiscalants inhibit sulfate and fluoride scales and disperse colloids and metal oxides. Despite claims that antiscalants can inhibit silica formation, no concrete evidence proves that silica polymerization is inhibited by antiscalants. Antiscalants can control acid-soluble scales at a fraction of the dosage required to control the same scale using sulfuric acid.
Some small-scale desalination units use 'beach wells'. These are usually drilled on the seashore. These intake facilities are relatively simple to build and the seawater they collect is pretreated via slow filtration through subsurface sand/seabed formations. Raw seawater collected using beach wells is often of better quality in terms of solids, silt, oil, grease, organic contamination, and microorganisms, compared to open seawater intakes. Beach intakes may also yield source water of lower salinity.
High pressure pump
The high pressure pump pushes water through the membrane. Typical pressures for brackish water range from 1.6 to 2.6 MPa (225 to 376 psi). In the case of seawater, they range from 5.5 to 8 MPa (800 to 1,180 psi). This requires substantial energy. Where energy recovery is used, part of the high pressure pump's work is done by the energy recovery device, reducing energy inputs.
Membrane assembly
The membrane assembly consists of a pressure vessel with a membrane that allows feedwater to be pushed against it. The membrane must be strong enough to withstand the pressure. RO membranes are made in a variety of configurations. The two most common are spiral-wound and hollow-fiber.
Only part of the water pumped onto the membrane passes through. The left-behind "concentrate" passes along the saline side of the membrane and flushes away the salt and other remnants. The percentage of desalinated water is the "recovery ratio". This varies with salinity and system design parameters: typically 20% for small seawater systems, 40% – 50% for larger seawater systems, and 80% – 85% for brackish water. The concentrate flow is typically 3 bar/50 psi less than the feed pressure, and thus retains much of the input energy.
The desalinated water purity is a function of the feed water salinity, membrane selection and recovery ratio. To achieve higher purity a second pass can be added which generally requires another pumping cycle. Purity expressed as total dissolved solids typically varies from 100 to 400 parts per million (ppm or mg/litre) on a seawater feed. A level of 500 ppm is generally the upper limit for drinking water, while the US Food and Drug Administration classifies mineral water as water containing at least 250 ppm.
Energy recovery
Energy recovery can reduce energy consumption by 50% or more. Much of the input energy can be recovered from the concentrate flow, and the increasing efficiency of energy recovery devices greatly reduces energy requirements. Devices used, in order of invention, are:
Turbine or Pelton wheel: a water turbine driven by the concentrate flow, connected to the pump drive shaft provides part of the input power. Positive displacement axial piston motors have been used in place of turbines on smaller systems.
Turbocharger: a water turbine driven by concentrate flow, directly connected to a centrifugal pump that boosts the output pressure, reducing the pressure needed from the pump and thereby its energy input, similar in construction principle to car engine turbochargers.
Pressure exchanger: using the pressurized concentrate flow, via direct contact or a piston, to pressurize part of the membrane feed flow to near concentrate flow pressure. A boost pump then raises this pressure by typically 3 bar / 50 psi to the membrane feed pressure. This reduces flow needed from the high-pressure pump by an amount equal to the concentrate flow, typically 60%, and thereby its energy input. These are widely used on larger low-energy systems. They are capable of 3 kWh/m3 or less energy consumption.
Energy-recovery pump: a reciprocating piston pump. The pressurized concentrate flow is applied to one side of each piston to help drive the membrane feed flow from the opposite side. These are the simplest energy recovery devices to apply, combining the high pressure pump and energy recovery in a single self-regulating unit. These are widely used on smaller low-energy systems. They are capable of 3 kWh/m3 or less energy consumption.
Batch operation: RO systems run with a fixed volume of fluid (thermodynamically a closed system) do not suffer from wasted energy in the brine stream, as the energy to pressurize a virtually incompressible fluid (water) is negligible. Such systems have the potential to reach second-law efficiencies of 60%. Such systems can be created multiple ways, including using pressurized tanks with pistons or bladders, or low-pressure tanks with conventional ERD's.
Remineralisation and pH adjustment
The desalinated water is stabilized to protect downstream pipelines and storage, usually by adding lime or caustic soda to prevent corrosion of concrete-lined surfaces. Liming material is used to adjust pH between 6.8 and 8.1 to meet the potable water specifications, primarily for effective disinfection and for corrosion control. Remineralisation may be needed to replace minerals removed from the water by desalination, although this process has proved to be costly and inconvenient in order to meet mineral demand by humans and plants as found in typical freshwater. For instance water from Israel's national water carrier typically contains dissolved magnesium levels of 20 to 25 mg/liter, while water from the Ashkelon plant has no magnesium. Ashkelon water created magnesium-deficiency symptoms in crops, including tomatoes, basil, and flowers, and had to be remedied by fertilization. Israeli drinking water standards require a minimum calcium level of 20 mg/liter. Askelon's post-desalination treatment uses sulfuric acid to dissolve calcite (limestone), resulting in calcium concentrations of 40 to 46 mg/liter, lower than the 45 to 60 mg/liter found in typical Israeli fresh water.
Disinfection
Post-treatment disinfection provides secondary protection against compromised membranes and downstream problems. Disinfection by means of ultraviolet (UV) lamps (sometimes called germicidal or bactericidal) may be employed to sterilize pathogens that evade the RO process. Chlorination or chloramination (chlorine and ammonia) protects against pathogens that may have lodged in the distribution system downstream.
Disadvantages
Large-scale industrial/municipal systems recover typically 75% to 80% of the feed water, or as high as 90%, because they can generate the required higher pressure.
Wastewater
Household RO units use a lot of water because they have low back pressure. Household RO water purifiers typically produce one liter of usable water and 3-25 liters of wastewater. The remainder is discharged, usually into the drain. Because wastewater carries the rejected contaminants, recovering this water is not practical for household systems. Wastewater is typically delivered to house drains. A RO unit delivering of treated water per day also discharge between . This led India's National Green Tribunal to propose a ban on RO water purification systems in areas where the total dissolved solids (TDS) measure in water is less than 500 mg/liter. In Delhi, large-scale use of household RO devices has increased the total water demand of the already water-parched National Capital Territory of India.
Health
RO removes both harmful contaminants and desirable minerals. Some studies report some relation between long-term health effects and consumption of water low on calcium and magnesium, although these studies are of low quality.
Waste-stream considerations
Depending upon the desired product, either the solvent or solute stream of RO will be waste. For food concentration applications, the concentrated solute stream is the product and the solvent stream is waste. For water treatment applications, the solvent stream is purified water and the solute stream is concentrated waste. The solvent waste stream from food processing may be used as reclaimed water, but there may be fewer options for disposal of a concentrated waste solute stream. Ships may use marine dumping and coastal desalination plants typically use marine outfalls. Landlocked RO plants may require evaporation ponds or injection wells to avoid polluting groundwater or surface runoff.
Research
Improving Current Membranes
Current RO membranes, thin-film composite (TFC) polyamide membranes, are being studied to find ways of improving their permeability. Through new imaging methods, researchers were able to make 3D models of membranes and examine how water flowed through them. They found that TFC membranes with areas of low flow significantly decreased water permeability. By ensuring uniformity of the membranes and allowing water to flow continuously without slowing down, membrane permeability could be improved by 30%-40%.
Electrodialysis
Research has examined integrating RO with electrodialysis to improve recovery of valuable deionized products, or to reduce concentrate volumes.
Low-pressure High-recovery (LPHR)
Another approach is low-pressure high-recovery multistage RO (LPHR). It produces concentrated brine and freshwater by cycling the output repeatedly through a relatively porous membrane at relatively low pressure. Each cycle removes additional impurities. Once the output is relatively pure, it is sent through a conventional RO membrane at conventional pressure to complete the filtration step. LPHR was found to be economically feasible, recovering more than 70% with an OPD between 58 and 65 bar and leaving no more than 350 ppm TDS from a seawater feed with 35,000 ppm TDS.
Carbon Nanotubes (CNTs)
Carbon nanotubes are meant to potentially solve the typical tradeoff between the permeability and the selectivity of RO membranes. CNTs present many ideal characteristics including: mechanical strength, electron affinity, and also exhibiting flexibility during modification. By restructuring carbon nanotubes and coating or impregnating them with other chemical compounds, scientists can manufacture these membranes to have all of the most desirable traits. The hope with CNT membranes is to find a combination of high water permeability while also decreasing the amount of neutral solutes taken out of the water. This would help decrease energy costs and the cost of remineralization after purification through the membrane.
Graphene
Graphene membranes are meant to take advantage of their thinness to increase efficiency. Graphene is a singular layer of carbon atoms, so it is about 1000 times thinner than existing membranes. Graphene membranes are around 100 nm thick while current membranes are about 100 μm. Many researchers were concerned with the durability of graphene and if it would be able to handle RO pressures. New research finds that depending on the substrate (a supporting layer that does no filtration and only provides structural support), graphene membranes can withstand 57MPa of pressure which is about 10 times the typical pressures for seawater RO.
Batch RO may offer increased energy efficiency, more durable equipment and higher salinity limits.
The conventional approach claimed that molecules cross the membrane individually. A research team devised a "solution-friction" theory, claiming that molecules in groups through transient pores. Characterizing that process could guide membrane development. The accepted theory is that individual water molecules diffuse through the membrane, termed the "solution-diffusion" model.
| Physical sciences | Other separations | Chemistry |
468974 | https://en.wikipedia.org/wiki/Fox%20Terrier | Fox Terrier | Fox Terriers are two different breeds of the terrier dog type: the Smooth Fox Terrier and the Wire Fox Terrier. Both of these breeds originated in the 19th century from a handful of dogs who are descended from earlier varieties of British terriers, and are related to other modern white terrier breeds. In addition, a number of breeds have diverged from these two main types of fox terrier and have been recognised separately, including the Jack Russell Terrier, Miniature Fox Terrier and Rat Terrier. The Wire and Smooth Fox Terriers share similar characteristics, the main differences being in the coat and markings. They have been successful in conformation shows, more prominently in America than their homeland.
History
English physician John Caius described the English terrier type in his 1577 work English Dogges. By the 18th century, it was recorded that all terriers were wire haired, and black and tan in colour. The earliest record of any white terrier was a dog named Pitch, who was owned by Colonel Thomas Thornton in 1790. The dog was the subject of a painting by Sawrey Gilpin, who created a portrait of it while it was still alive. An engraving of this painting, made prior to 1810, was accompanied by the text, "It would be necessary to notice Colonel Thorton's terriers if they were only on account of his justly celebrated Pitch, from whom are descended most of the white terriers in the kingdom." Dog writers of the early 20th century suggested that Pitch was a terrier-greyhound cross, which was how the colour was introduced into the breed.
Not much is known of early 19th century breeding practices that came to create the modern Fox Terrier. However it is thought that the Beagle, Old English Bulldog, English Toy Terrier, Pointer and even Dalmatian were all used in the creation and stabilisation of the breed.
From 1870 onwards, a complete pedigree for Fox Terriers exists. Three dogs, known as Old Jock, Trap and Tartar, are the ancestors of most modern strains of white terrier. Old Jock was bred from a black and tan terrier, while Trap was from the strain out of Reverend John "Jack" Russell's dogs. Russell would later have the Jack Russell Terrier series of breeds named after him. Although definitive history on Tartar is not known, he is also thought to have come from black and tan terrier stock. Of the female dogs from this period, Grove Nettle was the best known. In addition to those dogs, another named Old Tip is thought to be the forefather of the modern Wire Fox Terrier.
The terrier of this period was short of leg, built in a chunky fashion, and had a skull that was broad across the top. Used in hunting packs and kept in kennels along with English Foxhounds, they were bred for their ability to drive the fox out of its den. Gradually the pace of the hunt became faster, and the terriers were bred with longer legs for more speed towards the end of the 18th century. However the increased size meant that the Fox Terrier became too big to be used for its original purpose, and its use in fox hunting began to decline.
The Fox Terrier Club of England was founded in 1876, with the American Fox Terrier Club following nine years later, becoming the first breed speciality club to become a member of the American Kennel Club. A female Smooth Fox Terrier bred by Winthrop Rutherfurd named Warren Remedy was the first winner of Best in Show at the Westminster Kennel Club Dog Show in 1907, 1908, and 1909. Following the popularity of conformation dog shows, the Fox Terrier was the result of selective breeding. For instance, by 1949, 120 of 140 Wire Fox Terrier champions were descended from a single championship-winning dog.
Modern breeds
There are two main breeds of Fox Terrier, Smooth and Wire, both of which originate in England. In addition, there are several descendant breeds which have been developed in a variety of countries. For example, the American Toy Fox Terrier was developed from the Smooth Fox Terrier, Italian Greyhound, Manchester Terrier, Miniature Pinscher and Chihuahua breeds.
Both major types of Fox Terrier are mostly white with coloured markings and have similar facial characteristics. They are essentially the same except for differences in coats, and slight differences in colouration and markings. Markings on either type can appear black at birth on the head, but may lighten in adult life, the most common colours being tan and black. The markings are a type of piebald spotting. Although the two main breeds were originally commonly interbred, this no longer occurs in pure-bred lines.
The Fox Terrier has been more successful than any other breed at the Westminster Kennel Club Dog Show, with four Best in Show titles going to the Smooth Fox Terrier and fourteen titles going to the Wire Fox Terrier. They were both shown as one breed with two coat variations; this changed in 1985 when the American Kennel Club listed the two as separate breeds. In the UK, at Crufts, the Smooth Fox Terrier has not won any Best in Show titles, while the Wire Fox Terrier has won on three occasions.
Smooth Fox Terrier
The Smooth Fox Terrier has a short, hard coat which is predominantly white. It measures at the withers, and weighs between . The head of this breed is long and wedge shaped, with small, dark eyes and ears that are v-shaped.
The breed has been identified as one of several Vulnerable Native Breeds in the UK. This is when there are fewer than 300 annual registrations with The Kennel Club. In 2010, there were 155 Smooth Fox Terriers registered, compared to 693 for the Wire Fox Terrier and 8,663 for the most popular breed in the Terrier Group, the Staffordshire Bull Terrier.
The most successful dog at the Westminster Kennel Club Dog Show was Ch. Warren Remedy, who won the Best in Show title three times between 1907 and 1909. These victories were followed by a further Best in Show title for Ch. Sabine Rarebit in 1910. Despite the Smooth Fox Terrier breed winning the first four Best in Show titles at Westminster, it has not won again since.
Wire Fox Terrier
The Wire Haired Fox Terrier has a hard and crisp double coat with a coarse texture underneath that provides protection from the cold. It should be so dense that the skin cannot be seen or felt. The individual hairs should twist, but are not curly. An average member of the breed should measure at the withers and weigh . It has similar features to the Smooth Fox Terrier, with small dark eyes and v-shaped ears. Its body is shorter than it is tall.
King Edward VII owned a Wire Fox Terrier from the Notts kennel called Caesar of Notts, which did a great deal to popularise the breed at the turn of the 20th century. Another member of the Notts kennel was an early winner of Best Champion at Crufts in 1911 named Collarbone of Notts. Other individual dogs that greatly influenced the breed included Ch. Talavera Simon, born in 1924, and Ch. Zeloy Emperor, born in 1960.
Descendant breeds
Common health issues
Myasthenia gravis, a neuromuscular disease, is hereditary in the Smooth Fox Terrier through an autosomal recessive gene. This can also be a symptom of megaesophagus, which is a health issue for the Wire Fox Terrier. Another hereditary condition in the Smooth Fox Terrier is cataracts, which is more prevalent than average in the breed. Both types of Fox Terrier can be susceptible to allergies.
In a survey conducted by The Kennel Club, the primary cause of death for Fox Terriers was old age, causing 31.8% of reported deaths. The secondary cause was cancer of an unspecified type, which accounted for 22.7% of reports. The average lifespan of a Fox Terrier is around 15 years; the Kennel Club survey reported a median age at death of 13 years and 2 months.
The wire-haired Fox Terrier is predisposed to atopic dermatitis.
| Biology and health sciences | Dogs | Animals |
469220 | https://en.wikipedia.org/wiki/Frontal%20lobe | Frontal lobe | The frontal lobe is the largest of the four major lobes of the brain in mammals, and is located at the front of each cerebral hemisphere (in front of the parietal lobe and the temporal lobe). It is parted from the parietal lobe by a groove between tissues called the central sulcus and from the temporal lobe by a deeper groove called the lateral sulcus (Sylvian fissure). The most anterior rounded part of the frontal lobe (though not well-defined) is known as the frontal pole, one of the three poles of the cerebrum.
The frontal lobe is covered by the frontal cortex. The frontal cortex includes the premotor cortex and the primary motor cortex – parts of the motor cortex. The front part of the frontal cortex is covered by the prefrontal cortex. The nonprimary motor cortex is a functionally defined portion of the frontal lobe.
There are four principal gyri in the frontal lobe. The precentral gyrus is directly anterior to the central sulcus, running parallel to it and contains the primary motor cortex, which controls voluntary movements of specific body parts. Three horizontally arranged subsections of the frontal gyrus are the superior frontal gyrus, the middle frontal gyrus, and the inferior frontal gyrus. The inferior frontal gyrus is divided into three parts – the orbital part, the triangular part and the opercular part.
The frontal lobe contains most of the dopaminergic neurons in the cerebral cortex. The dopaminergic pathways are associated with reward, attention, short-term memory tasks, planning, and motivation. Dopamine tends to limit and select sensory information coming from the thalamus to the forebrain.
Structure
The frontal lobe is the largest lobe of the brain and makes up about a third of the surface area of each hemisphere. On the lateral surface of each hemisphere, the central sulcus separates the frontal lobe from the parietal lobe. The lateral sulcus separates the frontal lobe from the temporal lobe.
The frontal lobe can be divided into a lateral, polar, orbital (above the orbit; also called basal or ventral), and medial part. Each of these parts consists of a particular gyrus:
Lateral part: lateral part of the superior frontal gyrus, middle frontal gyrus, and inferior frontal gyrus
Polar part: frontopolar cortex, transverse frontopolar gyri, frontomarginal gyrus.
Orbital part: Lateral orbital gyrus, anterior orbital gyrus, posterior orbital gyrus, medial orbital gyrus, and gyrus rectus
Medial part: Medial part of the superior frontal gyrus, cingulate gyrus.
The gyri are separated by sulci. E.g., the precentral gyrus is in front of the central sulcus, and behind the precentral sulcus. The superior and middle frontal gyri are divided by the superior frontal sulcus. The middle and inferior frontal gyri are divided by the inferior frontal sulcus.
In humans the frontal lobe reaches full maturity only after the 20s—the prefrontal cortex, in particular, continues in maturing 'til the second and third decades of life—which, thereafter, marks the cognitive maturity associated with adulthood. A small amount of atrophy, however, is normal in the aging person's frontal lobe. Fjell, in 2009, studied atrophy of the brain in people aged 60–91 years. The 142 healthy participants were scanned using MRI. Their results were compared to those of 122 participants with Alzheimer's disease. A follow-up one year later showed there to have been a marked volumetric decline in those with Alzheimer's and a much smaller decline (averaging 0.5%) in the healthy group. These findings corroborate those of Coffey, who in 1992 indicated that the frontal lobe decreases in volume approximately 0.5–1% per year.
Function
The entirety of the frontal cortex can be considered the "action cortex", much as the posterior cortex is considered the "sensory cortex". It is devoted to action of one kind or another: skeletal movement, ocular movement, speech control, and the expression of emotions. In humans, the largest part of the frontal cortex, the prefrontal cortex (PFC), is responsible for internal, purposeful mental action, commonly called reasoning or prefrontal synthesis.
The function of the PFC involves the ability to project future consequences that result from current actions. PFC functions also include override and suppression of socially unacceptable responses as well as differentiation of tasks.
The PFC also plays an important part in integrating longer non-task based memories stored across the brain. These are often memories associated with emotions derived from input from the brain's limbic system. The frontal lobe modifies those emotions, generally to fit socially acceptable norms.
Psychological tests that measure frontal lobe function include finger tapping (as the frontal lobe controls voluntary movement), the Wisconsin Card Sorting Test, and measures of language, numeracy skills, and decision making, all of which are controlled by the frontal lobe.
Clinical significance
Damage
Damage to the frontal lobe can occur in a number of ways and result in many different consequences. Transient ischemic attacks (TIAs) also known as mini-strokes, and strokes are common causes of frontal lobe damage in older adults (65 and over). These strokes and mini-strokes can occur due to the blockage of blood flow to the brain or as a result of the rupturing of an aneurysm in a cerebral artery. Other ways in which injury can occur include traumatic brain injuries incurred following accidents, diagnoses such as Alzheimer's disease or Parkinson's disease (which cause dementia symptoms), and frontal lobe epilepsy (which can occur at any age). Very often, frontal lobe damage is recognized in those with prenatal alcohol exposure.
Symptoms
Common effects of damage to the frontal lobe are varied. Patients who have experienced frontal lobe trauma may know the appropriate response to a situation but display inappropriate responses to those same situations in real life. Similarly, emotions that are felt may not be expressed in the face or voice. For example, someone who is feeling happy would not smile, and the voice would be devoid of emotion. Along the same lines, though, the person may also exhibit excessive, unwarranted displays of emotion. Depression is common in stroke patients. Also common is a loss of or decrease in motivation. Someone might not want to carry out normal daily activities and would not feel "up to it". Those who are close to the person who has experienced the damage may notice changes in behavior. This personality change is characteristic of damage to the frontal lobe, and was exemplified in the case of Phineas Gage. The frontal lobe is the same part of the brain that is responsible for executive functions such as planning for the future, judgment, decision-making skills, attention span, and inhibition. These functions can decrease drastically in someone whose frontal lobe is damaged.
Consequences that are seen less frequently are also varied. Confabulation may be the most frequently indicated "less common" effect. In the case of confabulation, someone gives false information while maintaining the belief that it is the truth. In a small number of patients, uncharacteristic cheerfulness can be noted. This effect is seen mostly in patients with lesions to the right frontal portion of the brain.
Another infrequent effect is that of reduplicative paramnesia, in which patients believe that the location in which they currently reside is a replica of one located somewhere else. Similarly, those who experience Capgras syndrome after frontal lobe damage believe that an identical "replacement" has taken the identity of a close friend, relative, or other person and is posing as that person. This last effect is seen mostly in schizophrenic patients who also have a neurological disorder in the frontal lobe.
DNA damage
In the human frontal cortex, a set of genes undergo reduced expression after age 40 and especially after age 70. This set includes genes that have key functions in synaptic plasticity important in learning and memory, vesicular transport and mitochondrial function. During aging, DNA damage is markedly increased in the promoters of the genes displaying reduced expression in the frontal cortex. In cultured human neurons, these promoters are selectively damaged by oxidative stress.
Individuals with HIV associated neurocognitive disorders accumulate nuclear and mitochondrial DNA damage in the frontal cortex.
Genetic
A report from the National Institute of Mental Health says a gene variant of (COMT) that reduces dopamine activity in the prefrontal cortex is related to poorer performance and inefficient functioning of that brain region during working memory, tasks, and to a slightly increased risk for schizophrenia.
History
Psychosurgery
In the early 20th century, a medical treatment for mental illness, first developed by Portuguese neurologist Egas Moniz, involved damaging the pathways connecting the frontal lobe to the limbic system. A frontal lobotomy (sometimes called frontal leucotomy) successfully reduced distress but at the cost of often blunting the subject's emotions, volition and personality. The indiscriminate use of this psychosurgical procedure, combined with its severe side effects and a mortality rate of 7.4 to 17 per cent, earned it a bad reputation. The frontal lobotomy has largely died out as a psychiatric treatment. More precise psychosurgical procedures are still used, although rarely. They may include anterior capsulotomy (bilateral thermal lesions of the anterior limbs of the internal capsule) or the bilateral cingulotomy (involving lesions of the anterior cingulate gyri) and might be used to treat otherwise untreatable obsessional disorders or clinical depression.
Theories of function
Theories of frontal lobe function can be separated into four categories:
Single-process theories, which propose that "damage to a single process or system is responsible for a number of different dysexecutive symptoms"
Multi-process theories, which propose "that the frontal lobe executive system consists of a number of components that typically work together in everyday actions (heterogeneity of function)"
Construct-led theories, which propose that "most if not all frontal functions can be explained by one construct (homogeneity of function) such as working memory or inhibition"
Single-symptom theories, which propose that a specific dysexecutive symptom (e.g., confabulation) is related to the processes and construct of the underlying structures.
Other theories include:
Stuss (1999) suggests a differentiation into two categories according to homogeneity and heterogeneity of function.
Grafman's managerial knowledge units (MKU) / structured event complex (SEC) approach (cf. Wood & Grafman, 2003)
Miller & Cohen's integrative theory of prefrontal functioning (e.g. Miller & Cohen, 2001)
Rolls's stimulus-reward approach and Stuss's anterior attentional functions (Burgess & Simons, 2005; Burgess, 2003; Burke, 2007).
It may be highlighted that the theories described above differ in their focus on certain processes/systems or construct-lets. Stuss (1999) remarks that the question of homogeneity (single construct) or heterogeneity (multiple processes/systems) of function "may represent a problem of semantics and/or incomplete functional analysis rather than an unresolvable dichotomy" (p. 348). However, further research will show if a unified theory of frontal lobe function that fully accounts for the diversity of functions will be available.
Other primates
Many scientists had thought that the frontal lobe was disproportionately enlarged in humans compared to other primates. This was thought to be an important feature of human evolution and seen as the primary reason why human cognition differs from that of other primates. However, this view in relation to great apes has since been challenged by neuroimaging studies. Using magnetic resonance imaging to determine the volume of the frontal cortex in humans, all extant ape species, and several monkey species, it was found that the human frontal cortex was not relatively larger than the cortex of other great apes, but was relatively larger than the frontal cortex of lesser apes and the monkeys. The higher cognition of the humans is instead seen to relate to a greater connectedness given by neural tracts that do not affect the cortical volume. This is also evident in the pathways of the language network connecting the frontal and temporal lobes.
| Biology and health sciences | Nervous system | Biology |
469246 | https://en.wikipedia.org/wiki/Lattice%20%28group%29 | Lattice (group) | In geometry and group theory, a lattice in the real coordinate space is an infinite set of points in this space with the properties that coordinate-wise addition or subtraction of two points in the lattice produces another lattice point, that the lattice points are all separated by some minimum distance, and that every point in the space is within some maximum distance of a lattice point. Closure under addition and subtraction means that a lattice must be a subgroup of the additive group of the points in the space, and the requirements of minimum and maximum distance can be summarized by saying that a lattice is a Delone set. More abstractly, a lattice can be described as a free abelian group of dimension which spans the vector space . For any basis of , the subgroup of all linear combinations with integer coefficients of the basis vectors forms a lattice, and every lattice can be formed from a basis in this way. A lattice may be viewed as a regular tiling of a space by a primitive cell.
Lattices have many significant applications in pure mathematics, particularly in connection to Lie algebras, number theory and group theory. They also arise in applied mathematics in connection with coding theory, in percolation theory to study connectivity arising from small-scale interactions, cryptography because of conjectured computational hardness of several lattice problems, and are used in various ways in the physical sciences. For instance, in materials science and solid-state physics, a lattice is a synonym for the framework of a crystalline structure, a 3-dimensional array of regularly spaced points coinciding in special cases with the atom or molecule positions in a crystal. More generally, lattice models are studied in physics, often by the techniques of computational physics.
Symmetry considerations and examples
A lattice is the symmetry group of discrete translational symmetry in n directions. A pattern with this lattice of translational symmetry cannot have more, but may have less symmetry than the lattice itself. As a group (dropping its geometric structure) a lattice is a finitely-generated free abelian group, and thus isomorphic to .
A lattice in the sense of a 3-dimensional array of regularly spaced points coinciding with e.g. the atom or molecule positions in a crystal, or more generally, the orbit of a group action under translational symmetry, is a translation of the translation lattice: a coset, which need not contain the origin, and therefore need not be a lattice in the previous sense.
A simple example of a lattice in is the subgroup . More complicated examples include the E8 lattice, which is a lattice in , and the Leech lattice in . The period lattice in is central to the study of elliptic functions, developed in nineteenth century mathematics; it generalizes to higher dimensions in the theory of abelian functions. Lattices called root lattices are important in the theory of simple Lie algebras; for example, the E8 lattice is related to a Lie algebra that goes by the same name.
Dividing space according to a lattice
A lattice in thus has the form
where {v1, ..., vn} is a basis for . Different bases can generate the same lattice, but the absolute value of the determinant of the vectors vi is uniquely determined by and denoted by d(). If one thinks of a lattice as dividing the whole of into equal polyhedra (copies of an n-dimensional parallelepiped, known as the fundamental region of the lattice), then d() is equal to the n-dimensional volume of this polyhedron. This is why d() is sometimes called the covolume of the lattice. If this equals 1, the lattice is called unimodular.
Lattice points in convex sets
Minkowski's theorem relates the number d() and the volume of a symmetric convex set S to the number of lattice points contained in S. The number of lattice points contained in a polytope all of whose vertices are elements of the lattice is described by the polytope's Ehrhart polynomial. Formulas for some of the coefficients of this polynomial involve d() as well.
Computational lattice problems
Computational lattice problems have many applications in computer science. For example, the Lenstra–Lenstra–Lovász lattice basis reduction algorithm (LLL) has been used in the cryptanalysis of many public-key encryption schemes, and many lattice-based cryptographic schemes are known to be secure under the assumption that certain lattice problems are computationally difficult.
Lattices in two dimensions: detailed discussion
There are five 2D lattice types as given by the crystallographic restriction theorem. Below, the wallpaper group of the lattice is given in IUCr notation, Orbifold notation, and Coxeter notation, along with a wallpaper diagram showing the symmetry domains. Note that a pattern with this lattice of translational symmetry cannot have more, but may have less symmetry than the lattice itself. A full list of subgroups is available. For example, below the hexagonal/triangular lattice is given twice, with full 6-fold and a half 3-fold reflectional symmetry. If the symmetry group of a pattern contains an n-fold rotation then the lattice has n-fold symmetry for even n and 2n-fold for odd n.
For the classification of a given lattice, start with one point and take a nearest second point. For the third point, not on the same line, consider its distances to both points. Among the points for which the smaller of these two distances is least, choose a point for which the larger of the two is least. (Not logically equivalent but in the case of lattices giving the same result is just "Choose a point for which the larger of the two is least".)
The five cases correspond to the triangle being equilateral, right isosceles, right, isosceles, and scalene. In a rhombic lattice, the shortest distance may either be a diagonal or a side of the rhombus, i.e., the line segment connecting the first two points may or may not be one of the equal sides of the isosceles triangle. This depends on the smaller angle of the rhombus being less than 60° or between 60° and 90°.
The general case is known as a period lattice. If the vectors p and q generate the lattice, instead of p and q we can also take p and p-q, etc. In general in 2D, we can take a p + b q and c p + d q for integers a,b, c and d such that ad-bc is 1 or -1. This ensures that p and q themselves are integer linear combinations of the other two vectors. Each pair p, q defines a parallelogram, all with the same area, the magnitude of the cross product. One parallelogram fully defines the whole object. Without further symmetry, this parallelogram is a fundamental parallelogram.
The vectors p and q can be represented by complex numbers. Up to size and orientation, a pair can be represented by their quotient. Expressed geometrically: if two lattice points are 0 and 1, we consider the position of a third lattice point. Equivalence in the sense of generating the same lattice is represented by the modular group: represents choosing a different third point in the same grid, represents choosing a different side of the triangle as reference side 0–1, which in general implies changing the scaling of the lattice, and rotating it. Each "curved triangle" in the image contains for each 2D lattice shape one complex number, the grey area is a canonical representation, corresponding to the classification above, with 0 and 1 two lattice points that are closest to each other; duplication is avoided by including only half of the boundary. The rhombic lattices are represented by the points on its boundary, with the hexagonal lattice as vertex, and i for the square lattice. The rectangular lattices are at the imaginary axis, and the remaining area represents the parallelogrammatic lattices, with the mirror image of a parallelogram represented by the mirror image in the imaginary axis.
Lattices in three dimensions
The 14 lattice types in 3D are called Bravais lattices. They are characterized by their space group. 3D patterns with translational symmetry of a particular type cannot have more, but may have less, symmetry than the lattice itself.
Lattices in complex space
A lattice in is a discrete subgroup of which spans as a real vector space. As the dimension of as a real vector space is equal to , a lattice in will be a free abelian group of rank .
For example, the Gaussian integers form a lattice in , as is a basis of over .
In Lie groups
More generally, a lattice Γ in a Lie group G is a discrete subgroup, such that the quotient G/Γ is of finite measure, for the measure on it inherited from Haar measure on G (left-invariant, or right-invariant—the definition is independent of that choice). That will certainly be the case when G/Γ is compact, but that sufficient condition is not necessary, as is shown by the case of the modular group in SL2(R), which is a lattice but where the quotient isn't compact (it has cusps). There are general results stating the existence of lattices in Lie groups.
A lattice is said to be uniform or cocompact if G/Γ is compact; otherwise the lattice is called non-uniform.
Lattices in general vector spaces
While we normally consider lattices in this concept can be generalized to any finite-dimensional vector space over any field. This can be done as follows:
Let K be a field, let V be an n-dimensional K-vector space, let be a K-basis for V and let R be a ring contained within K. Then the R lattice in V generated by B is given by:
In general, different bases B will generate different lattices. However, if the transition matrix T between the bases is in - the general linear group of R (in simple terms this means that all the entries of T are in R and all the entries of are in R - which is equivalent to saying that the determinant of T is in - the unit group of elements in R with multiplicative inverses) then the lattices generated by these bases will be isomorphic since T induces an isomorphism between the two lattices.
Important cases of such lattices occur in number theory with K a p-adic field and R the p-adic integers.
For a vector space which is also an inner product space, the dual lattice can be concretely described by the set
or equivalently as
Related notions
A primitive element of a lattice is an element that is not a positive integer multiple of another element in the lattice.
| Mathematics | Other | null |
469304 | https://en.wikipedia.org/wiki/Lattice%20%28order%29 | Lattice (order) | A lattice is an abstract structure studied in the mathematical subdisciplines of order theory and abstract algebra. It consists of a partially ordered set in which every pair of elements has a unique supremum (also called a least upper bound or join) and a unique infimum (also called a greatest lower bound or meet). An example is given by the power set of a set, partially ordered by inclusion, for which the supremum is the union and the infimum is the intersection. Another example is given by the natural numbers, partially ordered by divisibility, for which the supremum is the least common multiple and the infimum is the greatest common divisor.
Lattices can also be characterized as algebraic structures satisfying certain axiomatic identities. Since the two definitions are equivalent, lattice theory draws on both order theory and universal algebra. Semilattices include lattices, which in turn include Heyting and Boolean algebras. These lattice-like structures all admit order-theoretic as well as algebraic descriptions.
The sub-field of abstract algebra that studies lattices is called lattice theory.
Definition
A lattice can be defined either order-theoretically as a partially ordered set, or as an algebraic structure.
As partially ordered set
A partially ordered set (poset) is called a lattice if it is both a join- and a meet-semilattice, i.e. each two-element subset has a join (i.e. least upper bound, denoted by ) and dually a meet (i.e. greatest lower bound, denoted by ). This definition makes and binary operations. Both operations are monotone with respect to the given order: and implies that and
It follows by an induction argument that every non-empty finite subset of a lattice has a least upper bound and a greatest lower bound. With additional assumptions, further conclusions may be possible; see Completeness (order theory) for more discussion of this subject. That article also discusses how one may rephrase the above definition in terms of the existence of suitable Galois connections between related partially ordered sets—an approach of special interest for the category theoretic approach to lattices, and for formal concept analysis.
Given a subset of a lattice, meet and join restrict to partial functions – they are undefined if their value is not in the subset The resulting structure on is called a . In addition to this extrinsic definition as a subset of some other algebraic structure (a lattice), a partial lattice can also be intrinsically defined as a set with two partial binary operations satisfying certain axioms.
As algebraic structure
A lattice is an algebraic structure , consisting of a set and two binary, commutative and associative operations and on satisfying the following axiomatic identities for all elements (sometimes called ):
The following two identities are also usually regarded as axioms, even though they follow from the two absorption laws taken together. These are called .
These axioms assert that both and are semilattices. The absorption laws, the only axioms above in which both meet and join appear, distinguish a lattice from an arbitrary pair of semilattice structures and assure that the two semilattices interact appropriately. In particular, each semilattice is the dual of the other. The absorption laws can be viewed as a requirement that the meet and join semilattices define the same partial order.
Connection between the two definitions
An order-theoretic lattice gives rise to the two binary operations and Since the commutative, associative and absorption laws can easily be verified for these operations, they make into a lattice in the algebraic sense.
The converse is also true. Given an algebraically defined lattice one can define a partial order on by setting
for all elements The laws of absorption ensure that both definitions are equivalent:
and dually for the other direction.
One can now check that the relation introduced in this way defines a partial ordering within which binary meets and joins are given through the original operations and
Since the two definitions of a lattice are equivalent, one may freely invoke aspects of either definition in any way that suits the purpose at hand.
Bounded lattice
A bounded lattice is a lattice that additionally has a (also called , or element, and denoted by or and a (also called , or , denoted by or by which satisfy
A bounded lattice may also be defined as an algebraic structure of the form such that is a lattice, (the lattice's bottom) is the identity element for the join operation and (the lattice's top) is the identity element for the meet operation
A partially ordered set is a bounded lattice if and only if every finite set of elements (including the empty set) has a join and a meet. For every element of a poset it is vacuously true that
and
and therefore every element of a poset is both an upper bound and a lower bound of the empty set. This implies that the join of an empty set is the least element and the meet of the empty set is the greatest element This is consistent with the associativity and commutativity of meet and join: the join of a union of finite sets is equal to the join of the joins of the sets, and dually, the meet of a union of finite sets is equal to the meet of the meets of the sets, that is, for finite subsets and of a poset
and
hold. Taking to be the empty set,
and
which is consistent with the fact that
Every lattice can be embedded into a bounded lattice by adding a greatest and a least element. Furthermore, every non-empty finite lattice is bounded, by taking the join (respectively, meet) of all elements, denoted by (respectively ) where is the set of all elements.
Connection to other algebraic structures
Lattices have some connections to the family of group-like algebraic structures. Because meet and join both commute and associate, a lattice can be viewed as consisting of two commutative semigroups having the same domain. For a bounded lattice, these semigroups are in fact commutative monoids. The absorption law is the only defining identity that is peculiar to lattice theory. A bounded lattice can also be thought of as a commutative without the distributive axiom.
By commutativity, associativity and idempotence one can think of join and meet as operations on non-empty finite sets, rather than on pairs of elements. In a bounded lattice the join and meet of the empty set can also be defined (as and respectively). This makes bounded lattices somewhat more natural than general lattices, and many authors require all lattices to be bounded.
The algebraic interpretation of lattices plays an essential role in universal algebra.
Examples
For any set the collection of all subsets of (called the power set of ) can be ordered via subset inclusion to obtain a lattice bounded by itself and the empty set. In this lattice, the supremum is provided by set union and the infimum is provided by set intersection (see Pic. 1).
For any set the collection of all finite subsets of ordered by inclusion, is also a lattice, and will be bounded if and only if is finite.
For any set the collection of all partitions of ordered by refinement, is a lattice (see Pic. 3).
The positive integers in their usual order form an unbounded lattice, under the operations of "min" and "max". 1 is bottom; there is no top (see Pic. 4).
The Cartesian square of the natural numbers, ordered so that if The pair is the bottom element; there is no top (see Pic. 5).
The natural numbers also form a lattice under the operations of taking the greatest common divisor and least common multiple, with divisibility as the order relation: if divides is bottom; is top. Pic. 2 shows a finite sublattice.
Every complete lattice (also see below) is a (rather specific) bounded lattice. This class gives rise to a broad range of practical examples.
The set of compact elements of an arithmetic complete lattice is a lattice with a least element, where the lattice operations are given by restricting the respective operations of the arithmetic lattice. This is the specific property that distinguishes arithmetic lattices from algebraic lattices, for which the compacts only form a join-semilattice. Both of these classes of complete lattices are studied in domain theory.
Further examples of lattices are given for each of the additional properties discussed below.
Examples of non-lattices
Most partially ordered sets are not lattices, including the following.
A discrete poset, meaning a poset such that implies is a lattice if and only if it has at most one element. In particular the two-element discrete poset is not a lattice.
Although the set partially ordered by divisibility is a lattice, the set so ordered is not a lattice because the pair 2, 3 lacks a join; similarly, 2, 3 lacks a meet in
The set partially ordered by divisibility is not a lattice. Every pair of elements has an upper bound and a lower bound, but the pair 2, 3 has three upper bounds, namely 12, 18, and 36, none of which is the least of those three under divisibility (12 and 18 do not divide each other). Likewise the pair 12, 18 has three lower bounds, namely 1, 2, and 3, none of which is the greatest of those three under divisibility (2 and 3 do not divide each other).
Morphisms of lattices
The appropriate notion of a morphism between two lattices flows easily from the above algebraic definition. Given two lattices and a lattice homomorphism from L to M is a function such that for all
Thus is a homomorphism of the two underlying semilattices. When lattices with more structure are considered, the morphisms should "respect" the extra structure, too. In particular, a bounded-lattice homomorphism (usually called just "lattice homomorphism") between two bounded lattices and should also have the following property:
In the order-theoretic formulation, these conditions just state that a homomorphism of lattices is a function preserving binary meets and joins. For bounded lattices, preservation of least and greatest elements is just preservation of join and meet of the empty set.
Any homomorphism of lattices is necessarily monotone with respect to the associated ordering relation; see Limit preserving function. The converse is not true: monotonicity by no means implies the required preservation of meets and joins (see Pic. 9), although an order-preserving bijection is a homomorphism if its inverse is also order-preserving.
Given the standard definition of isomorphisms as invertible morphisms, a is just a bijective lattice homomorphism. Similarly, a is a lattice homomorphism from a lattice to itself, and a is a bijective lattice endomorphism. Lattices and their homomorphisms form a category.
Let and be two lattices with 0 and 1. A homomorphism from to is called 0,1-separating if and only if ( separates 0) and ( separates 1).
Sublattices
A of a lattice is a subset of that is a lattice with the same meet and join operations as That is, if is a lattice and is a subset of such that for every pair of elements both and are in then is a sublattice of
A sublattice of a lattice is a of if and implies that belongs to for all elements
Properties of lattices
We now introduce a number of important properties that lead to interesting special classes of lattices. One, boundedness, has already been discussed.
Completeness
A poset is called a if its subsets have both a join and a meet. In particular, every complete lattice is a bounded lattice. While bounded lattice homomorphisms in general preserve only finite joins and meets, complete lattice homomorphisms are required to preserve arbitrary joins and meets.
Every poset that is a complete semilattice is also a complete lattice. Related to this result is the interesting phenomenon that there are various competing notions of homomorphism for this class of posets, depending on whether they are seen as complete lattices, complete join-semilattices, complete meet-semilattices, or as join-complete or meet-complete lattices.
"Partial lattice" is not the opposite of "complete lattice" – rather, "partial lattice", "lattice", and "complete lattice" are increasingly restrictive definitions.
Conditional completeness
A conditionally complete lattice is a lattice in which every subset has a join (that is, a least upper bound). Such lattices provide the most direct generalization of the completeness axiom of the real numbers. A conditionally complete lattice is either a complete lattice, or a complete lattice without its maximum element its minimum element or both.
Distributivity
Since lattices come with two binary operations, it is natural to ask whether one of them distributes over the other, that is, whether one or the other of the following dual laws holds for every three elements :
Distributivity of over
Distributivity of over
A lattice that satisfies the first or, equivalently (as it turns out), the second axiom, is called a distributive lattice.
The only non-distributive lattices with fewer than 6 elements are called M3 and N5; they are shown in Pictures 10 and 11, respectively. A lattice is distributive if and only if it does not have a sublattice isomorphic to M3 or N5. Each distributive lattice is isomorphic to a lattice of sets (with union and intersection as join and meet, respectively).
For an overview of stronger notions of distributivity that are appropriate for complete lattices and that are used to define more special classes of lattices such as frames and completely distributive lattices, see distributivity in order theory.
Modularity
For some applications the distributivity condition is too strong, and the following weaker property is often useful. A lattice is if, for all elements the following identity holds:
()
This condition is equivalent to the following axiom:
implies ()
A lattice is modular if and only if it does not have a sublattice isomorphic to N5 (shown in Pic. 11).
Besides distributive lattices, examples of modular lattices are the lattice of submodules of a module (hence modular), the lattice of two-sided ideals of a ring, and the lattice of normal subgroups of a group. The set of first-order terms with the ordering "is more specific than" is a non-modular lattice used in automated reasoning.
Semimodularity
A finite lattice is modular if and only if it is both upper and lower semimodular. For a graded lattice, (upper) semimodularity is equivalent to the following condition on the rank function
Another equivalent (for graded lattices) condition is Birkhoff's condition:
for each and in if and both cover then covers both and
A lattice is called lower semimodular if its dual is semimodular. For finite lattices this means that the previous conditions hold with and exchanged, "covers" exchanged with "is covered by", and inequalities reversed.
Continuity and algebraicity
In domain theory, it is natural to seek to approximate the elements in a partial order by "much simpler" elements. This leads to the class of continuous posets, consisting of posets where every element can be obtained as the supremum of a directed set of elements that are way-below the element. If one can additionally restrict these to the compact elements of a poset for obtaining these directed sets, then the poset is even algebraic. Both concepts can be applied to lattices as follows:
A continuous lattice is a complete lattice that is continuous as a poset.
An algebraic lattice is a complete lattice that is algebraic as a poset.
Both of these classes have interesting properties. For example, continuous lattices can be characterized as algebraic structures (with infinitary operations) satisfying certain identities. While such a characterization is not known for algebraic lattices, they can be described "syntactically" via Scott information systems.
Complements and pseudo-complements
Let be a bounded lattice with greatest element 1 and least element 0. Two elements and of are complements of each other if and only if:
In general, some elements of a bounded lattice might not have a complement, and others might have more than one complement. For example, the set with its usual ordering is a bounded lattice, and does not have a complement. In the bounded lattice N5, the element has two complements, viz. and (see Pic. 11). A bounded lattice for which every element has a complement is called a complemented lattice.
A complemented lattice that is also distributive is a Boolean algebra. For a distributive lattice, the complement of when it exists, is unique.
In the case that the complement is unique, we write and equivalently, The corresponding unary operation over called complementation, introduces an analogue of logical negation into lattice theory.
Heyting algebras are an example of distributive lattices where some members might be lacking complements. Every element of a Heyting algebra has, on the other hand, a pseudo-complement, also denoted The pseudo-complement is the greatest element such that If the pseudo-complement of every element of a Heyting algebra is in fact a complement, then the Heyting algebra is in fact a Boolean algebra.
Jordan–Dedekind chain condition
A chain from to is a set where
The length of this chain is n, or one less than its number of elements. A chain is maximal if covers for all
If for any pair, and where all maximal chains from to have the same length, then the lattice is said to satisfy the Jordan–Dedekind chain condition.
Graded/ranked
A lattice is called graded, sometimes ranked (but see Ranked poset for an alternative meaning), if it can be equipped with a rank function sometimes to , compatible with the ordering (so whenever ) such that whenever covers then The value of the rank function for a lattice element is called its rank.
A lattice element is said to cover another element if but there does not exist a such that
Here, means and
Free lattices
Any set may be used to generate the free semilattice The free semilattice is defined to consist of all of the finite subsets of with the semilattice operation given by ordinary set union. The free semilattice has the universal property. For the free lattice over a set Whitman gave a construction based on polynomials over s members.
Important lattice-theoretic notions
We now define some order-theoretic notions of importance to lattice theory. In the following, let be an element of some lattice is called:
Join irreducible if implies for all If has a bottom element some authors require . When the first condition is generalized to arbitrary joins is called completely join irreducible (or -irreducible). The dual notion is meet irreducibility (-irreducible). For example, in Pic. 2, the elements 2, 3, 4, and 5 are join irreducible, while 12, 15, 20, and 30 are meet irreducible. Depending on definition, the bottom element 1 and top element 60 may or may not be considered join irreducible and meet irreducible, respectively. In the lattice of real numbers with the usual order, each element is join irreducible, but none is completely join irreducible.
Join prime if implies Again some authors require , although this is unusual. This too can be generalized to obtain the notion completely join prime. The dual notion is meet prime. Every join-prime element is also join irreducible, and every meet-prime element is also meet irreducible. The converse holds if is distributive.
Let have a bottom element 0. An element of is an atom if and there exists no element such that Then is called:
Atomic if for every nonzero element of there exists an atom of such that
Atomistic if every element of is a supremum of atoms.
However, many sources and mathematical communities use the term "atomic" to mean "atomistic" as defined above.
The notions of ideals and the dual notion of filters refer to particular kinds of subsets of a partially ordered set, and are therefore important for lattice theory. Details can be found in the respective entries.
| Mathematics | Order theory | null |
469343 | https://en.wikipedia.org/wiki/Mandrill | Mandrill | The mandrill (Mandrillus sphinx) is a large Old World monkey native to west central Africa. It is one of the most colorful mammals in the world, with red and blue skin on its face and posterior. The species is sexually dimorphic, as males have a larger body, longer canine teeth and brighter coloring. It is the largest monkey in the world. Its closest living relative is the drill, with which it shares the genus Mandrillus. Both species were traditionally thought to be baboons, but further evidence has shown that they are more closely related to white-eyelid mangabeys.
Mandrills mainly live in tropical rainforests but will also travel across savannas. They are active during the day and spend most of their time on the ground. Their preferred foods are fruit and seeds, but mandrills will consume leaves, piths, mushrooms, and animals from insects to juvenile bay duiker. Mandrills live in large, stable groups known as "hordes" which can number in the hundreds. Females form the core of these groups, while adult males are solitary and only reunite with the larger groups during the breeding season. Dominant males have the most vibrant colors and fattest flanks and rumps, and have the most success siring young.
The mandrill is classified as vulnerable on the IUCN Red List. Its biggest threats are habitat destruction and hunting for bushmeat. Gabon is considered the stronghold for the species. Its habitat has declined in Cameroon and Equatorial Guinea, while its range in the Republic of the Congo is limited.
Etymology
The word mandrill is derived from the English words man and drill—the latter meaning or and being West African in origin—and dated to 1744. The name appears to have originally referred to chimpanzees. The first scholar to record the name for the colorful monkey was Georges-Louis Buffon in 1766. It was called the "tufted ape", "great baboon" and "ribbernosed baboon" by Thomas Pennant in A Synopsis of Quadrupeds (1771) and A History of Quadrupeds (1781).
Taxonomy
The mandrill was first scientifically depicted in Historia animalium (1551–1558) by Conrad Gessner, who considered it a kind of hyena. The species was formally classified by Carl Linnaeus as Simia sphinx in 1758. Its current generic name Mandrillus was coined by Ferdinand Ritgen in 1824.
Historically, some scientists placed the mandrill and the closely related drill (M. leucophaeus) in the baboon genus Papio. Morphological and genetic studies in the late 20th and early 21st centuries found a closer relationship to white-eyelid mangabeys of the genus Cercocebus. Some have even proposed that the mandrill and drill belong to Cercocebus. Two genetic studies in 2011 clarified Mandrillus and Cercocebus as separate sister lineages. The two genera split around 4.5 million years ago (mya) while the mandrill and drill split approximately 3.17 mya. Fossils of Mandrillus have not been found.
Some authorities have divided mandrill populations into subspecies: the northern mandrill (M. s. sphinx) and the southern mandrill (M. s. madarogaster). A proposed third subspecies, M. s. insularis, was based on the mistaken belief that mandrills are present on Bioko Island. The consensus is that mandrills belong to one subspecies (M. s. sphinx).
Cytochrome-b sequences suggest that mandrill populations north and south of the Ogooué River split 800,000 years ago and belong to distinct haplogroups. This divergence appears to have also led to the splitting of the mandrill strain of the simian immunodeficiency virus (SIV). The draft (incomplete) genome of the mandrill was published in 2020, with a reported genome size of 2.90 giga–base-pairs and high levels of heterozygosity.
Characteristics
The mandrill has a stocky body with a large head and muzzle, as well as a short and stumpy tail. The limbs are evenly sized and the fingers and toes are more elongated than those in baboons, with a more opposable big toe on the feet. The mandrill is the most sexually dimorphic primate, and it is the largest monkey. Females are less stocky and have shorter, flatter snouts. Males have a head-body length and weigh while females have a head-body length and weigh . Most of the teeth are larger in males, and the canine teeth reach up to and long for males and females respectively. Both sexes have long tails.
The coat of the mandrill is primarily grizzled or banded olive-brown with a yellow-orange beard and sparse, light hairs on its underside. The lips are surrounded by stiff white whiskers, and white bare skin exists behind the ears. Male mandrills have a "crest" of long hairs on the head and neck, while both sexes have chest glands which are covered by long hairs. The face, rump and genitals have less hair. The mandrill has a red line running down the middle of their face which connects to their red nose. On either side of the line, the skin is blue and grooved. In males, the blue skin is supported by ridged bone swellings. Females have more subdued facial coloring, but this can vary between individuals with some having stronger red and blue hues and others being darker or almost black. In males, the rump and areas around the genitals are multi-colored, consisting of red, pink, blue and purple skin, with a red penis shaft and violet scrotum. The genital and anal areas of the female are red.
The mandrill is among the most colorful mammals. Charles Darwin wrote in The Descent of Man: "no other member of the whole class of mammals is coloured in so extraordinary a manner as the adult male mandrill". The red coloration is created by blood vessels near the surface of the skin, while the blue is a form of structural coloration caused by parallel arrangements of collagen fibers. The blue ridges on males contrast with both the red facial hues and the green foliage of their environment, helping them stand out to other individuals. The darker and more subdued coloring of female faces is caused by melanin.
Ecology
The mandrill lives in west-central Africa, including southern Cameroon, mainland Equatorial Guinea (Río Muni), Gabon and parts of the Republic of the Congo. Its range is bounded by the Sanaga River to the north and the Ogooué and Ivindo Rivers to the east. It does not appear to share habitat with the drill, as the two species are separated by the Sanaga River. Mandrills live in tropical rainforests, generally preferring primary forests over secondary forests. They also live in patchy gallery forests surrounded by savanna and travel across grass areas within their forest habitats. They have also been recorded in mountainous areas, near rivers and in cultivated fields.
Mandrills prefer thick bush dominated by perennial plants like gingers and plants of the genera Brillantaisia and Phaulopsis. They mainly dwell on the ground, but feed as high as the canopy. Both mandrills and drills are more arboreal than baboons. Mandrills may aggregate or compete with other primates such as talapoins, guenons, mangabeys, black-and-white colobuses, chimpanzees, and gorillas.
Feeding
The mandrill is an omnivore. The core of its diet consists of plants, of which it eats over a hundred species. One study found the mandrill's diet was composed of fruit (50.7%), seeds (26.0%), leaves (8.2%), pith (6.8%), flowers (2.7%), and animal matter (4.1%), with other foods making up the remaining 1.4%. During the wet season, mandrills forage in continuous forest, when fruit is most available, while during the dry season they feed in gallery forests and at the borders of savannas and forests.
The mandrill's preferred fruits include those of the cashew species Pseudospondias microcarpa, the coffee species Nauclea diderrichii and the wort species Psorospermum febrifugum. Mandrills consume more seeds than many other primate species. Adult male mandrills are one of the few primates capable of biting through the hard shell of Detarium microcarpum seeds. For vegetation, they mostly eat the young leaves, shoots and piths of monocot plants. In particular, mandrills consume leaves from the arrowroots Haumania liebrechtsiana and Trachyphrynium braunianum, as well as the piths of ginger plants like Renealmia macrocolia and species in the genus Aframomum. They are also known to consume mushrooms.
The rest of a mandrill's diet is largely made up of invertebrates, particularly ants, termites, crickets, spiders, snails, and scorpions. They also eat birds and their eggs, frogs and rodents. Mandrills have been recorded preying on larger vertebrates such as juvenile bay duikers. Such prey is killed with a bite to the head followed by pulling off the hind limbs and tearing open the belly. Individuals may cooperate during hunting and share kills.
Predators, parasites and pathogens
Leopards may prey on mandrills, as traces of mandrill have been found in their feces. Other potential predators include African rock pythons, crowned eagles and chimpanzees. Leopards are a threat to all individuals, while eagles are only threats to the young. In a study where a mandrill group was exposed to models of leopards and crown eagles, the leopard models tended to cause the mandrills to flee up trees while the eagles were more likely to drive them to take cover. The dominant male did not flee from either model types; in the case of the leopards, he paced around while looking in their direction. Alarm calls were more commonly heard in response to leopards than eagles.
Mandrills can become infected with gastrointestinal parasites, such as nematodes and protozoa. Tumbu fly larvae may live under the skin and individuals that walk though grassland can get infested with ticks. Blood parasites include the malaria-causing Plasmodium and the nematode Loa loa, which is transmitted by bites from deer flies. Wild mandrills have tested positive for SIV, enteroviruses of the species EV-J and astroviruses, including a human variant.
Behavior and life history
Mandrills are mostly diurnal and are awake around 10 hours per day from morning to dusk. They often pick a new tree to sleep in every night. Mandrills have been observed using tools; in captivity, they use sticks to clean themselves. In the wild, mandrills appear to live 12–14 years, but captive individuals can live 30–40 years.
Social structure
Mandrills live in large "supergroups" or "hordes" that can contain hundreds of individuals. These large groups are fairly stable and do not appear to be gatherings of smaller ones. At Lopé National Park, Gabon, mandrill hordes were found to have an average of 620 individuals, and some groups were as large as 845, making them possibly the largest cohesive groups of wild primates. Another study in Lopé found that a horde of 625 mandrills consisted of 21 dominant males, 71 less dominant and subadult males, 247 adult and adolescent females, 200 juveniles, and 86 dependent infants. A mandrill horde of around 700 individuals in northern Lopé had a total home range of , of which was suitable habitat. The supergroup would occasionally diverge into two to four subgroups before reuniting. Another 15-month long study of a 120-member group found a home range of with an average traveling distance of per day.
Hordes consist of matrilineal family groups, and females are important for maintaining social cohesion. Strong connections with their relatives may lead to support during conflicts, higher survival rate of offspring and a longer lifespan for females. Dominant females are at the center of the group network and their removal leads to fewer social connections in the group. The social rank of a mother mandrill can contribute to the social rank of both her female and male offspring. Mature males are not permanent members of hordes but join as females become sexually receptive and leave as their sexual cycle ends. As a result, the coloration of the male mandrill may be intended to attract attention in a social structure with no long-term relationships between mates. Higher ranking males are found in the center of a social group while lower ranking males are more likely to occupy the periphery. Females have some control over the males and coalitions can expel an unwanted male from a group. Outside the breeding season, males are believed to lead a solitary life and all-male bachelor groups are not known to exist.
Both male and female mandrills rub and mark trees and branches with secretions from their chest glands, though males (and especially dominant males) mark more than females. The chemicals in the secretions signal the individual's sex, age and rank. Scent-marking may also serve a territorial function, captive alpha males will mark enclosure boundaries. Mandrills will groom one another, even when there is no benefit to be gained from doing so. During grooming, subordinates prefer to pick at other mandrills from behind, in order to minimize eye contact and give them more time to flee if the more dominant individual attacks. The recipients of grooming will try to maneuver the groomer to pick at more "risky" areas.
Reproduction and development
Dominant or alpha male mandrills have the most mating success. Upon gaining alpha status, males develop larger testicles, redder faces and posteriors, more secretion from the chest glands and fatter sides and rumps. When a male loses dominance, these physiological changes are at least partially reversed. The blue facial skin is more consistent in brightness. Higher ranking males tend to have more contrast between red and blue facial coloring. Due to their distribution of fat, dominant males are also known as "fatted" males while subordinate males are known as "non-fatted" males. Canine length also correlates with dominance, and males are less likely to sire offspring when their canines are under . In some individuals, the development of secondary sexual characteristics is suppressed in response to competition from other males. Male mandrills tend to establish dominance with vocalizations and facial expressions, rather than fighting.
Mating occurs mostly during the dry season, with female ovulation peaking between June and September. Receptive females have sexual swellings on their posteriors, and the red facial coloration can communicate age and fertility. Males also appear to detect a female's reproductive state using the vomeronasal organ (known as the flehmen response). Dominant males try to monopolize access to females by mate guarding, which involves the male tending to and copulating with a female for days. Dominant males tend to sire most of the offspring, but they are less able to monopolize access to the females when many females reach estrus at the same time. A subordinate male is also more likely to have reproductive success if he is closely related to an alpha male. An ovulating female tends to allow the brightest colored males to come near her and touch her perineum, and is more likely to groom and solicit them. The female signals her willingness to mate by positioning her posterior towards the male. Intercourse lasts no more than 60 seconds, with the male mounting the female and making pelvic thrusts.
Mandrill gestation lasts an average of 175 days with most births taking place between January and March, during the wet season. Gaps in between births range from 184 to 1,159 days with an average of 405 days. and tend to be shorter in higher ranking females. Infants are born at an average weight of , and mostly bare-skinned with some white hair and a tuft of dark hair on the head and along the spine. Over the next two or three months, they develop their adult hair color on the body, limbs and head while the flesh-colored face and snout darken. Dependent infants are carried on their mothers' bellies. Young are typically weaned at around 230 days old. Males become more sexually dimorphic between four and eight years old, at which point females are already beginning to give birth. Males start leaving their horde after they reach six years old. Females reach their adult size around seven years while males do so at ten years.
Communication
Mandrills communicate with various facial expressions and postures. Threat displays involve open mouth staring, usually in combination with head bobbing, ground slapping and raised hair. These gestures are usually performed by dominant individuals towards subordinates, who respond with bared teeth grimaces, signaling fear and aggression. Both young and low-ranking females show submission and anxiety with a pouting "duck face". Playful intentions are communicated with a relaxed open-mouth face. Males approaching females display a "grin" or silent bared-teeth face and make lip-smacks. This display may also occur with teeth-chattering. Mandrills can develop and pass on new gestures; captive individuals at the Colchester Zoo, England facepalm to discourage being disturbed, particularly while resting.
Mandrills also produce several vocalizations, for both long and short distances. During group movements, adult males produce two-phase grunts and one-syllable roars, both of which are equivalent to the "wahoo" bark of baboons. Other group members produce "crowings", which last almost two seconds and start as a vibration and transition into a longer harmonic sound. Short distance vocals include the "yak", a sharp, repeating, pulse-like call produced by all individuals except for adult males and made in tense situations. Mandrills may also grunt during aggressive encounters. Growls are used to express mild alarm while intense alarms come in the form of a short, two-syllable sharp call known as the "k-alarm". A sharp, loud "K-sound" is produced for unknown reasons. Screaming is a signal of fear and made by individuals fleeing, while the girney, a type of moan or purr, is made as a form of appeasement or frustration among females and young. Individual voices are more similar among related animals, but unrelated mandrills can have similar voices if they regularly interact.
Threats and conservation
As of 2019, the IUCN Red List lists the mandrill as vulnerable. Its total population is unknown but is suspected to have decreased by more than 30 percent over the last 24 years. Its main threats are habitat destruction and hunting for bushmeat. The mandrill appears to have suffered massive habitat loss in Equatorial Guinea and southern Cameroon, while its range in the Republic of the Congo is limited and its status is unknown. In addition, while mandrills live in groups numbering in the hundreds, hunting in Cameroon and Equatorial Guinea appears to have led to smaller group sizes. Gabon is seen as the most important remaining refuge for the species, and the country's low population density and vast rainforests make it a good candidate for mandrill conservation. Surveys have shown high population numbers for other primate species like chimpanzees and gorillas. A semi-wild population exists at the International Centre of Medical Research of Franceville.
The mandrill is listed under Appendix I by CITES, banning commercial trade in wild-caught specimens, and under Class B by the African Convention, which provides them protection but allows special authorization for their killing, capturing or collecting. There is at least one protected area for mandrills within each of the countries they inhabit. In Gabon, most of the rainforests have been leased to timber companies but around 10 percent is part of a national parks system, 13 of which were established in 2002.
| Biology and health sciences | Old World monkeys | Animals |
469646 | https://en.wikipedia.org/wiki/Watercress | Watercress | Watercress or yellowcress (Nasturtium officinale) is a species of aquatic flowering plant in the cabbage family, Brassicaceae.
Watercress is a rapidly growing perennial plant native to Eurasia. It is one of the oldest known leaf vegetables consumed by humans. Watercress and many of its relatives, such as garden cress, mustard, radish, and wasabi, are noteworthy for their piquant flavors.
Description
Watercress can grow up to in length. The stems are hollow and float in water. The leaf structure is pinnately compound. Small, white, and green inflorescences are produced in clusters and are frequently visited by insects, especially hoverflies, such as Eristalis flies.
Taxonomy
Watercress is listed in some sources as belonging to the genus Rorippa, although molecular evidence shows those aquatic species with hollow stems are more closely related to Cardamine than Rorippa. Despite the Latin name, watercress is not particularly closely related to the flowers popularly known as nasturtiums (Tropaeolum majus). T. majus belongs to the family Tropaeolaceae, a sister taxon to the Brassicaceae within the order Brassicales.
Distribution and habitat
In some regions, watercress is regarded as a weed, in other regions as an aquatic vegetable or herb. Watercress has grown in many temperate locations worldwide.
Watercress was introduced into China through Hong Kong and Macao in the 1800s from Europe. Its cultivation also spread to highland areas in the tropical regions of Asia.
Clear fast-flowing chalk streams are the primary natural habitat for wild watercress in the United Kingdom. Many settlements in England are named after watercress, from Old English êacerse, including Kersey, Kesgrave, Kersal, and Kershopefoot.
Health concerns
Watercress crops grown in the presence of manure can be an environment for parasites such as the liver fluke, Fasciola hepatica. Cultivated watercress has the advantage of being free of the liver fluke.
When introduced into non-native environments watercress can have negative impacts on native species. With the introduction of watercress, the organic matter in the sediment increases which in turn attracts predatory macroinvertebrates that feed on other plants in the environment.
By inhibiting the cytochrome P450 enzyme CYP2E1, compounds in watercress may alter drug metabolism in individuals on certain medications such as chlorzoxazone.
Due to its fast-growing nature and invasive species status, Nasturtium officinale is prohibited in Illinois.
Uses
Watercress leaves, stems, and fruit can be eaten raw. In China watercress is often boiled alongside pork and traditional medicinal ingredients to make a wintertime tonic soup, in Vietnam it is generally used raw as a component in salads.
Tradition
Ancient Romans thought eating it would cure mental illness. Twelfth-century mystic Hildegard of Bingen thought eating it steamed and drinking the water would cure jaundice or fever. Watercress was eaten by Native Americans. Some Native Americans used it to treat kidney illnesses and constipation, and it was thought by some to be an aphrodisiac. Early African Americans used the plant as an abortifacient; it was believed to cause sterility as well.
Nutrition
The new tips of watercress leaves can be eaten raw or cooked, although caution should be used when collecting these in the wild because of parasites such as giardia. Watercress is 95% water and has low contents of carbohydrates, protein, fat, and dietary fiber. A 100-gram serving of raw watercress provides , is particularly rich in vitamin K (238% of the Daily Value, DV), and contains significant amounts of vitamin A, vitamin C, riboflavin, vitamin B6, calcium, and manganese.
Phytochemicals and cooking
As a cruciferous vegetable, watercress contains isothiocyanates that are partly destroyed by boiling, while the bioavailability of its carotenoids is slightly increased by cooking. Steaming or microwave cooking retains these phytochemicals a bit better than boiling.
Cultivation
Watercress cultivation is practical on both a large scale and a garden scale. Being semi-aquatic, watercress is well-suited to hydroponic cultivation, thriving best in water that is slightly alkaline. It is frequently produced around the headwaters of chalk streams. In many local markets, the demand for hydroponically grown watercress exceeds supply, partly because cress leaves are unsuitable for distribution in dried form and can only be stored fresh for about 2–3 days.
Also sold as sprouts, the edible shoots are harvested days after germination. If unharvested, watercress can grow to a height of .
In the United Kingdom, watercress was first commercially cultivated in 1808 by the horticulturist William Bradbery along the River Ebbsfleet in Kent. Historically important areas of cultivation also included Hampshire, Stamford, and Watercress Wildlife Site at St Albans. Watercress is now grown in several counties, most notably Hampshire, Dorset, Wiltshire, and Hertfordshire. New Alresford in Hampshire is considered to be the nation's watercress capital, and its railway, the Watercress Line, was named for this cargo.
| Biology and health sciences | Leafy vegetables | Plants |
469661 | https://en.wikipedia.org/wiki/Whiteboard | Whiteboard | A whiteboard (also known by marker board, dry-erase board, dry-wipe board, and pen-board) is a glossy, usually white surface for making non-permanent markings. Whiteboards are analogous to blackboards, but with a smoother surface allowing for rapid marking and erasing of markings on their surface. The popularity of whiteboards increased rapidly in the mid-1990s and they have become a fixture in many offices, meeting rooms, school classrooms, public events and other work environments.
The term whiteboard is also used metaphorically in reference to features of computer software applications that simulate whiteboards. Such "virtual tech whiteboards" allow one or more people to write or draw images on a simulated canvas. This is a common feature of many virtual meetings, collaborations, and instant messaging applications. The term whiteboard is also used to refer to interactive whiteboards.
History
It has been widely reported that Korean War veteran and photographer Martin Heit and Albert Stallion, an employee at Alliance, a leading flat rolled steel sheet supplier should be credited with the invention of the whiteboard in the 1950s. Heit and Stallion may have popularized and/or perfected whiteboards, but in reality the history of whiteboards long precedes them. A thorough examination of the invention of whiteboards reveals the concept was introduced two decades earlier by mechanical engineer Paul F. Born who installed one in a classroom in Elgin, Ill., in 1937. Born was the head of the district's school board at the time.
Whiteboards became commercially available in the early 1960s, but did not become widely used until 30 years later. Early whiteboards needed to be wiped with a damp cloth and markers had a tendency to leave marks behind, even after the board was erased.
In 1974, whiteboards were proposed as additional equipment for Soviet schools.
Dry-erase markers for whiteboards were invented in 1975.
Whiteboards began being commonly used by businesses in the early 1990s. They became more common in classrooms during the 1990s due to concerns over health problems in children with dust allergies and the potential for chalk dust to damage computers. By the late 1990s, about 21% of American classrooms had converted from chalkboards to whiteboards.
Types
The first whiteboards were very expensive and were made of an enameled steel. Cheaper versions were then produced, including laminated chipboard, high-pressure laminates, and steel boards with a matte-finished or glossy white, usually polyester or acrylic coating. Enameled whiteboards, also referred to as porcelain and sometimes glass boards, have the advantage that markings can be erased completely; other materials tend to become stained over time. Enameled boards are more expensive and less used in commercial environments, but in more demanding environments with heavier use, such as educational establishments, porcelain boards are considered superior.
Other types of dry marker boards are also available, such as high gloss vinyl and coated paper, which can be rolled up, high-density two-part high gloss paints, glass and coated acrylics and polypropylene magic whiteboards which use static electricity to cling to walls, windows, and doors.
Clear marker surfaces, made of glass (matte or glossy) or specially coated acrylic, became available around the mid-2000s. They are generally manufactured from technical glass, e.g. for monitor screen filters, which is optically coated.
Adhesive whiteboards
Whiteboard material can be bought in rolls, sheets, and pre-formed boards. Adhesive whiteboards come in either a sheet or a roll and feature a stick back enabling the user to create a custom size board or project with the material. Although adhesive whiteboard material does not come in a thick, hard glass or painted steel plate, the melamine allows for a flexible material while preserving the high-quality whiteboard attributes of other surface materials. Adhesive whiteboards allow for custom projects such as dry erase wall calendars, whiteboard tables, cupboard grocery lists, indoor games for kids, and more.
Erasable marker
One type of whiteboard pen (also called a whiteboard marker or dry erase marker) was invented by Jerry Woolf of Techform Laboratories and later patented by Pilot Pen in 1975. It is a non-permanent marker and uses an erasable ink that adheres to the writing surface without binding to or being absorbed by it. Applications range from temporary writing with acetate sheets (for use with overhead projectors) to whiteboards and similar glossy surfaces. The erasable ink does not contain the toxic chemical compounds xylene and/or toluene, unlike permanent markers.
Surface materials
There are six types of materials commonly used for whiteboard surfaces:
Melamine
A resin-infused paper which is typically used over a substrate that can range from particle board to MDF (medium density fiberboard). Melamine boards range in quality primarily because of the amount of resin deposited on the base material. Some melamine boards remain clean (no ghosting) for a long time, others less so. Generally, this least expensive type of whiteboard is most commonly found in use in non-institutional applications. They are available in any office supply store. Performance varies widely. These boards are not suitable for heavy use, as in many educational cases, with time the paint erodes and the original surface reappears.
Steel or aluminum
Painted steel and aluminum dry erase also have a wide range of quality. Painted surfaces tend to be smoother, which leads to better methods of erasing. The painted surface is generally a multiple layer of coatings made up of a base coat in color (most commonly white) and a clear performance coating that is the dry erase component. Paint varies from electron beam cured coatings to UV and other coating systems. Good commercial grade painted steel or aluminum has excellent dry erase properties and many will be able to have permanent marker cleaned from the surface. Any coated surface is susceptible to scratching. Painted steel surfaces are magnetic and allow the use of magnets. Painted aluminum surfaces are rarely used as a base for whiteboards as they are not magnetic and are more expensive than steel. Painted steel whiteboards are most commonly used for custom printed whiteboards. These products are used as tracking boards, patient information boards and tournament and training boards.
Hardcoat laminate
Here again the performance varies over a wide range depending on the amount of resin used in the manufacturer. This category primarily uses melamine as its dry erase performance coat. Depending on the manufacturer (and the price) these laminates often are less porous and highly resistant to staining. Less common than other whiteboard surfaces, because they usually are used in combination with something else (a cabinet, doors or table tops for example). (This statement applies to Porcelain steel which is the only lifetime warranty available on the market)
Porcelain, enamel-on-steel
Ceramic fired onto a steel surface in a kiln. They are the most durable surfaces, and most carry a lifetime warranty. They are very common in heavy-use industrial settings. They are highly scratch-resistant, although materials harder than glass, such as diamond, can scratch them. They do not absorb dry erase or permanent marker ink. They allow the use of magnets. The surface can be cleaned with any non-abrasive cleaner suitable for porcelain, then rinsed off with water to prevent smearing. Permanent marker can be removed using a polar solvent such as ethanol, isopropanol or acetone, or by writing over it with a dry-erase marker and erasing it.
Tempered glass
Tempered glass is easier to erase. Most tempered glass white boards have no frame.
Polypropylene film
Polypropylene film is temporary and uses static electricity to cling to existing surfaces. It works with dry whiteboard markers and wipes. Whiteboards are reusable and materials are recyclable.
| Technology | Writing tools | null |
469990 | https://en.wikipedia.org/wiki/Whirlpool | Whirlpool | A whirlpool is a body of rotating water produced by opposing currents or a current running into an obstacle. Small whirlpools form when a bath or a sink is draining. More powerful ones formed in seas or oceans may be called maelstroms ( ). Vortex is the proper term for a whirlpool that has a downdraft.
In narrow ocean straits with fast flowing water, whirlpools are often caused by tides. Many stories tell of ships being sucked into a maelstrom, although only smaller craft are actually in danger. Smaller whirlpools appear at river rapids and can be observed downstream of artificial structures such as weirs and dams. Large cataracts, such as Niagara Falls, produce strong whirlpools.
Notable whirlpools
Saltstraumen
Saltstraumen is a narrow strait located close to the Arctic Circle, south-east of the city of Bodø, Norway.
It has one of the strongest tidal currents in the world. Whirlpools up to in diameter and in depth are formed when the current is at its strongest.
Moskstraumen
Moskstraumen or Moske-stroom is an unusual system of whirlpools in the open seas in the Lofoten Islands off the Norwegian coast. It is the second strongest whirlpool in the world with flow currents reaching speeds as high as . This is supposedly the whirlpool depicted in Olaus Magnus's map, labeled as "Horrenda Caribdis" (Charybdis).
The Moskstraumen is formed by the combination of powerful semi-diurnal tides and the unusual shape of the seabed, with a shallow ridge between the Moskenesøya and Værøya islands which amplifies and whirls the tidal currents.
The fictional depictions of the Moskstraumen by Edgar Allan Poe, Jules Verne, and Cixin Liu describe it as a gigantic circular vortex that reaches the bottom of the ocean, when in fact it is a set of currents and crosscurrents with a rate of . Poe described this phenomenon in his short story "A Descent into the Maelström", which in 1841 was the first to use the word maelstrom in the English language; in this story related to the Lofoten Maelstrom, two fishermen are swallowed by the maelstrom while one survives.
Corryvreckan
The Corryvreckan is a narrow strait between the islands of Jura and Scarba, in Argyll and Bute, on the northern side of the Gulf of Corryvreckan, Scotland. It is the third-largest whirlpool in the world. Flood tides and inflow from the Firth of Lorne to the west can drive the waters of Corryvreckan to waves of more than , and the roar of the resulting maelstrom, which reaches speeds of , can be heard away. Though it was classified initially as non-navigable by the Royal Navy it was later categorized as "extremely dangerous".
A documentary team from Scottish independent producers Northlight Productions once threw a mannequin into the Corryvreckan ("the Hag") with a high-visibility vest and depth gauge. The mannequin was swallowed and spat up far down current with a depth gauge reading of and evidence of being dragged along the bottom for a great distance.
Niagara Whirlpool
About three miles (4.8 kilometers) downstream from Niagara Falls is the Niagara Whirlpool. Located mostly in Canada and partially in the United States, the whirlpool is crossed by the Whirlpool Aero Car.
The basin of the whirlpool is 1,700 feet (518 meters) long and 1,200 feet (365 meters) wide. Its maximum water depth is 125 feet (38 meters).
Other notable maelstroms and whirlpools
Old Sow whirlpool is located between Deer Island, New Brunswick, Canada, and Moose Island, Eastport, Maine, USA. It is given the epithet "pig-like" as it makes a screeching noise when the vortex is at its full fury and reaches speeds of as much as . The smaller whirlpools around this Old Sow are known as "Piglets".
The Naruto whirlpools are located in the Naruto Strait near Awaji Island in Japan, which have speeds of .
Skookumchuck Narrows is a tidal rapids that develops whirlpools, on the Sunshine Coast, British Columbia, Canada with speeds of the current exceeding .
French Pass () is a narrow and treacherous stretch of water that separates D'Urville Island from the north end of the South Island of New Zealand. In 2000 a whirlpool there caught student divers, resulting in fatalities.
A short-lived whirlpool sucked in a portion of the Lake Peigneur in Louisiana, United States after a drilling mishap on November 20, 1980. This was not a naturally occurring whirlpool, but a disaster caused by underwater drillers breaking through the roof of a salt mine. The lake then drained into the mine until the mine filled and the water levels equalized, but the formerly deep lake was now deep. This mishap caused a sinkhole, and in the end, resulted in the destruction of five houses, the loss of nineteen barges and eight tug boats, oil rigs, a mobile home, trees, acres of land, and most of a botanical garden. The adjacent settlement of Jefferson Island was reduced in area by 10%. A crater across was left behind. Nine of the barges, which had sunk, later resurfaced after the whirlpool subsided.
A more recent example of an artificial whirlpool that received significant media coverage occurred in early June 2015, when an intake vortex formed in Lake Texoma, on the Oklahoma–Texas border, near the floodgates of the dam that forms the lake. At the time of the whirlpool's formation, the lake was being drained after reaching its highest level ever. The Army Corps of Engineers, which operates the dam and lake, expected that the whirlpool would last until the lake reached normal seasonal levels by late July.
Dangers
Powerful whirlpools have killed unlucky seafarers, but their power tends to be exaggerated by laymen. One of the few reports of a large disaster comes from the fourteenth-century Mali Empire ruler Mansa Musa, as reported by a contemporary, Ibn Fadlallah al-Umari:
Tales like those by Paul the Deacon, Edgar Allan Poe, and Jules Verne are entirely fictional.
However, temporary whirlpools caused by major engineering disasters, such as the Lake Peigneur disaster, have been recorded as capable of submerging medium-sized watercraft such as barges and tugboats.
In literature and popular culture
Besides Poe and Verne, another literary source is of the 1500s, Olaus Magnus, a Swedish bishop, who had stated that a maelstrom more powerful than the one written about in the Odyssey sucked in ships, which sank to the bottom of the sea, and even whales were pulled in. Pytheas, the Greek historian, also mentioned that maelstroms swallowed ships and threw them up again.
The monster Charybdis of Greek mythology was later rationalized as a whirlpool, which sucked entire ships into its fold in the narrow coast of Sicily, a disaster faced by navigators.
During the 8th century, Paul the Deacon, who had lived among the Belgii, described tidal bores and the maelstrom for a Mediterranean audience unused to such violent tidal surges:
Three of the most notable literary references to the Lofoten Maelstrom date from the nineteenth century. The first is a short story by Edgar Allan Poe named "A Descent into the Maelström" (1841). The second is Twenty Thousand Leagues Under the Seas (1870), a novel by Jules Verne. At the end of this novel, Captain Nemo seems to commit suicide, sending his Nautilus submarine into the Maelstrom (although in Verne's sequel Nemo and the Nautilus were seen to have survived). The "Norway maelstrom" is also mentioned in Herman Melville's Moby-Dick.
In the Life of St Columba, the author, Adomnan of Iona, attributes to the saint miraculous knowledge of a particular bishop who sailed into a whirlpool off the coast of Ireland. In Adomnan's narrative, he quotes Columba saying
The Corryvreckan whirlpool plays a key role in the 1945 Powell and Pressburger film I Know Where I'm Going!. Joan Webster (Wendy Hiller) is determined to get to the Isle of Kiloran and marry her fiancé. Dangerous weather delays her crossing, and her determination becomes desperate when she realizes that she is falling in love with Torquil MacNeil (Roger Livesey). Against the advice of experienced folk, she offers a young fisherman a huge sum of money to take her over. At the last moment, Royal Naval Officer Torquil steps into the boat, and after a squall knocks the engine out of commission, they face the whirlpool. Torquil manages to repair the engine before the tide turns, and they return to the mainland. This part of the picture uses footage Powell filmed, while tied to a mast to leave both hands free for the camera, at Corryvreckan, incorporated into scenes shot in a huge tank at the studio.
In the 2007 film Pirates of the Caribbean: At World's End, the final battle between the Black Pearl and the Flying Dutchman takes place with both ships sailing inside a giant whirlpool which appears to be over a kilometer wide and several hundred meters deep. The fantasy novels Eldest and The Bellmaker (otherwise unconnected) both feature a scene where the protagonists' ship escapes pursuit by successfully navigating a massive whirlpool, while the pursuing vessel fails to do so and is dragged under.
Etymology
One of the earliest uses in English of the Scandinavian word malström or malstrøm was by Edgar Allan Poe in his short story "A Descent into the Maelström" (1841). The Nordic word itself is derived from the Dutch word maelstrom (; modern spelling ), from malen ('to mill' or 'to grind') and stroom ('stream'), to form the meaning 'grinding current' or literally 'mill-stream', in the sense of milling (grinding) grain.
| Physical sciences | Water: General | Earth science |
470323 | https://en.wikipedia.org/wiki/Iron%28II%29%20oxide | Iron(II) oxide | Iron(II) oxide or ferrous oxide is the inorganic compound with the formula FeO. Its mineral form is known as wüstite. One of several iron oxides, it is a black-colored powder that is sometimes confused with rust, the latter of which consists of hydrated iron(III) oxide (ferric oxide). Iron(II) oxide also refers to a family of related non-stoichiometric compounds, which are typically iron deficient with compositions ranging from Fe0.84O to Fe0.95O.
Preparation
FeO can be prepared by the thermal decomposition of iron(II) oxalate.
FeC2O4 → FeO + CO2 + CO
The procedure is conducted under an inert atmosphere to avoid the formation of iron(III) oxide (Fe2O3). A similar procedure can also be used for the synthesis of manganous oxide and stannous oxide.
Stoichiometric FeO can be prepared by heating Fe0.95O with metallic iron at 770 °C and 36 kbar.
Reactions
FeO is thermodynamically unstable below 575 °C, tending to disproportionate to metal and Fe3O4:
4FeO → Fe + Fe3O4
Structure
Iron(II) oxide adopts the cubic, rock salt structure, where iron atoms are octahedrally coordinated by oxygen atoms and the oxygen atoms octahedrally coordinated by iron atoms. The non-stoichiometry occurs because of the ease of oxidation of FeII to FeIII effectively replacing a small portion of FeII with two-thirds their number of FeIII, which take up tetrahedral positions in the close packed oxide lattice.
In contrast to the crystalline solid, in the molten state iron atoms are coordinated by predominantly 4 or 5 oxygen atoms.
Below 200 K there is a minor change to the structure which changes the symmetry to rhombohedral and samples become antiferromagnetic.
Occurrence in nature
Iron(II) oxide makes up approximately 9% of the Earth's mantle. Within the mantle, it may be electrically conductive, which is a possible explanation for perturbations in Earth's rotation not accounted for by accepted models of the mantle's properties.
Uses
Iron(II) oxide is used as a pigment. It is FDA-approved for use in cosmetics and it is used in some tattoo inks. It can also be used as a phosphate remover from home aquaria.
| Physical sciences | Oxide salts | Chemistry |
470488 | https://en.wikipedia.org/wiki/Pencil%20sharpener | Pencil sharpener | A pencil sharpener (or pencil pointer, or in Ireland a parer or topper) is a tool for sharpening a pencil's writing point by shaving away its worn surface. Pencil sharpeners may be operated manually or by an electric motor. It is common for many sharpeners to have a casing around them, which can be removed for emptying the pencil shavings debris into a bin.
History
Before the development of dedicated pencil sharpeners, a pencil was sharpened by whittling with a knife.
The development of pencil sharpeners began in France when a French book from 1822 reported in detail about an invention of Mr. C. A. Boucher (Paris) for the construction of a pencil sharpener. He was working with pantographs and apparently needed a device to precisely sharpen the pencils. The device of Mr. Boucher was technically sensible and functional. His idea was also internationally known and recognized, as shown by corresponding reports in German literature at this time. But Mr. Boucher had not applied a patent for his pencil sharpener. Commercial use of his inventions is unlikely.
French mathematician Bernard Lassimonne (Limoges) applied for the world's first patent (French patent #2444) on a pencil sharpener in 1828. Pencil sharpener devices using his patent were actually produced and sold by Binant, a shop for painting accessories in Paris. In 1833 in England, Cooper & Eckstein patented the so-called Styloxynon, a simple device consisting of two sharp files set together at right angle in a small block of rosewood. This is the oldest pencil sharpener that has surviving examples.
In the 1830s and 1840s, some French people, all based in Paris, were engaged in construction of simple pencil sharpening tools, like François Joseph Lahausse. These devices were partially sold, but without supra-regional significance. In 1847, the French nobleman Thierry des Estivaux invented a simple hand-held pencil sharpener in its recognizable modern form. The first American pencil sharpener was patented by Walter Kittredge Foster of Bangor, Maine in 1855. He founded a company – the first pencil sharpener company in the world – and produced such small hand-held pencil sharpeners in a large amount. Only a few years later the sharpeners were sold also in Europe as "American pencil sharpeners".
At the end of the 19th century, especially in the United States, pencil sharpeners with various mechanisms had been developed and put on the market. These devices were often heavy and intended for use in offices. Examples are the Perfect Pencil Pointer (Goodell. Co.), the GEM Pencil Sharpener (by Gould & Cook Co.), the Planetary Pencil Sharpener (A. B. Dick Co.), all from the US or the Jupiter (Guhl & Harbeck Co.) from Germany. At the beginning of the 20th century the company Automatic Pencil Sharpener Co. (APSCO) was founded and brought out the US Automatic Pencil Sharpener after 1907, which dominated in those years. They later sold machines with milling mechanisms, such as the Climax, Dexter, Wizard, and Junior models. In the next few decades, APSCO became the largest pencil sharpening machine producer in the world and together with a few other US companies, it dominated the market.
Electric pencil sharpeners for offices have been made since at least 1917. A school teacher, Neil Carruthers, brought electric sharpeners from his visit to USA back to his town of Whitehaven for use in schools at the end of the Victorian Era.
In May 2011, tourism officials in Logan, Ohio put on display, in its regional welcome center, hundreds of pencil sharpeners which had been collected by Rev. Paul Johnson, an Ohio minister who died in 2010. Johnson, a World War II veteran, had kept his collection of more than 3,400 sharpeners in a small shed, outside his home in Carbon Hill in southeast Ohio. He had started collecting after his wife gave him a few pencil sharpeners as a gift in the late 1980s and kept them organized into categories, including cats, Christmas, and Disneyland.
Manual sharpeners
Prism sharpeners
So-called "prism" sharpeners, also called "manual" or "pocket" sharpeners in the United States, have no separate moving parts and are typically the smallest and cheapest commonly used pencil sharpener on the market. The simplest common variety is a small rectangular prism or block, only about 1 × 5/8 × 7/16 inch (2.5 × 1.7 × 1.1 cm) in size. The block-shaped sharpener consists of a combined point-shaping cone that is aligned to the cylindrical pencil alignment guide hole, into which the pencil is inserted. A sharp blade is mounted so that its sharp edge just enters the shaping cone tangentially. The pencil is inserted into the sharpener and rotated while the sharpener is held motionless. The body of the sharpener is often contoured, ridged or grooved to make the small block easier to firmly grip, and is typically made of aluminum alloy, magnesium alloy or hard plastic.
The blade inside the sharpener shaves the wood and graphite tip of the pencil, while the shavings emerge through a slot along the blade edge. It is important that the cylindrical alignment hole closely fits the diameter of the pencil, to keep the pencil from wobbling, which would cause stepped or lurching cut-depths and point breakage. Another important feature is a larger clearance hole at the end of the cone allowing sections of the pencil lead which break away to be removed with only minor inconvenience. Prism sharpeners can be bare or enclosed in a container to collect the shavings, while some enclosed sharpeners may be harder to clear in the event of a blockage.
A few prism sharpeners are hand-cranked, rotating the cutting blade instead of rotating the pencil. Moderate care is needed to not break the tip of the pencil being sharpened, requiring the pencil to be sharpened again. However, because pencils may have different standard diameters in different nations, imported sharpeners may have non-standard-sized alignment guide-holes, making sharpening attempts difficult. If the alignment hole is too small, the pencil cannot be inserted, while if it is too large, the tip of the pencil will repeatedly break off. Prism sharpeners may be right- or left-handed, requiring clockwise or counter-clockwise rotation of the pencil being sharpened.
Linear blade sharpeners
Unlike prism sharpeners, linear blade sharpeners do not rotate relative to the pencil being sharpened, and may be viewed as just a special form of knife, with a mechanical guide for increased safety and convenience. Some models use replaceable shaving razor blades, while others have permanently-fitted blades. Linear blade sharpeners may require more skill, but they allow one to sharpen the tip of the pencil into any desired shape and angle of taper, whereas prism sharpeners have a fixed sharpening angle and produce circular symmetry.
While most linear blade sharpeners are simple and directly hand-operated, some devices in the past were crank-operated, using mechanisms to convert crank rotation into linear motion.
Cylindrical (planetary) sharpeners
These mechanisms are also called planetary sharpeners, in reference to their use of planetary gears. A larger, stationary planetary sharpener can be mounted on a desk or wall and powered by a hand crank. Typically, the pencil is inserted into the sharpener with one hand, and the crank is turned with the other. This rotates a set of helical cylindrical cutters in the mechanism, set at an acute angle to each other. The multiple cutting edges quickly sharpen the pencil, with a more precise finish than a single-blade device. Some cylindrical sharpeners have only one helical cutter cylinder, but most have two cylinders or more.
Most planetary sharpeners have a large opening, with a rotatable guide disk in front of it that has multiple holes of different sizes, to accommodate pencils of many different diameters. Advanced models have a spring-driven holder for the pencil, so that the pencil automatically is pushed into the mechanism while being sharpened. Some versions also offer a regulator of the desired sharpness, since it is not always desired to make the graphite core needle-sharp.
Other systems
Some older models like the 1897 German Jupiter 1 used reversible rotary cutter-disks with cutting edges radiating from the center on each side. These were high-end models, quite large and expensive. Others simply used abrasives like sandpaper. In some cases an abrasive was used to shape the graphite core, while the wood was cut some other way.
Electric sharpeners
The oldest surviving electric pencil sharpener is the Boston Polar Club pencil sharpener, introduced around 1936. Electric pencil sharpeners work on the same principle as manual ones, but one or more flat-bladed or cylindrical cutters are rotated by an electric motor. Some electric pencil sharpeners are powered by batteries rather than being plugged into a building's electrical system, making them more portable.
Auto-stop electric pencil sharpeners are able to sense when the tip of the pencil is long enough, so they stop automatically. In basic automatic pencil sharpeners, the lead may become too long and break, and so users must be careful to supervise the operation.
Specialized pencil sharpeners
Specialized sharpeners are available that operate on non-standard sizes of pencil-shaped markers, such as wax crayons used in primary schools. Sharpeners that have two openings, one for normal pencils and one for larger crayons, are fairly common. Sharpeners with a single blade for use on wax crayons are available, and sometimes included in boxes of crayons. These often have plastic blades, which are adequate for the soft wax.
An artist's or draftsman's pencil sharpener leaves the graphite untouched and sharpens only the wood (some models can switch from standard to wood-only by an adjustment). The graphite lead is then honed to a sharp point with a lead pointer, which sharpens only the lead without wood. Lead pointers are also used with mechanical leadholders, with thicker diameter leads like 2 mm which have removable/refillable leads. Some sharpeners which function as a long point sharpener, have a second hole in which the blade sharpens the untouched graphite to a long, more precise point than would be otherwise possible using a single hole long point sharpener.
Carpenters may use carpenter pencils, the flattened shape of which stops them from rolling away, while still providing a constant line width. These pencils were traditionally sharpened with tools conveniently to hand, such as a plane or sandpaper. Rotating pencil sharpeners are now available for these pencils, in which a rotating plastic collar holds the pencil in position, although they then sharpen to the usual conical point as for a round pencil, abandoning some distinctive aspects of the carpenter's pencil. Alternatively, a special carpenter's pencil sharpener can be used, which has a sliding mechanism that leaves flat facets on the lead, in a manner similar to hand sharpening with a sharp knife.
Mechanical pencils with thin diameter leads dispense the graphite lead progressively during use and thus do not require sharpening; such pencils are sometimes called "self-sharpening". A type of mechanical pencil has a rotating gear mechanism which rotates the lead slightly every time the lead is lifted off the paper, helping to maintain a consistent, sharp point. If a finer or broader line is needed, a separate mechanical pencil using a lead with a different diameter is required.
Gallery
| Technology | Writing tools | null |
470605 | https://en.wikipedia.org/wiki/Red%20junglefowl | Red junglefowl | The red junglefowl (Gallus gallus), also known as the Indian red junglefowl (and formerly the bankiva or bankiva-fowl), is a species of tropical, predominantly terrestrial bird in the fowl and pheasant family, Phasianidae, found across much of Southeast and parts of South Asia. The red junglefowl was the primary species to give rise to today's many breeds of domesticated chicken (G. g. domesticus); additionally, the related grey junglefowl (G. sonneratii), Sri Lankan junglefowl (G. lafayettii) and the Javanese green junglefowl (G. varius) have also contributed genetic material to the gene pool of the modern chicken.
Molecular evidence, derived from whole-genome sequencing, has revealed that the chicken was first domesticated from red junglefowl ca. 8,000 years ago, with this domestication-event involving multiple maternal origins. Since then, the domestic form has spread around the world, and they are bred by humans in their millions for meat, eggs, colourful plumage and companionship. Outside of their native range, mainly in the Americas and Europe, the wild form of Gallus gallus is sometimes used in zoos, parks or botanical gardens as a free-ranging form of beneficial "pest control", similarly to—and often kept with—the Indian blue peafowl (Pavo cristatus) or the helmeted guineafowl (Numida meleagris); hybridisation has been documented between Gallus and Numida.
Taxonomy and systematics
The red jungle fowl has 5 recognized subspecies:
G. g. bankiva (Temminck, 1813) - Java and Bali
G. g. gallus (Linnaeus, 1758) - southern Myanmar through Indochina
G. g. jabouillei (Delacour & Kinnear, 1928) - south China to northern Vietnam and northern Laos
G. g. murghi (Robinson & Kloss, 1920) - north India, Nepal, Bhutan, and Bangladesh
G. g. spadiceus (Bonnaterre, 1792) - northeastern India to south China, the Malay peninsula, and north Sumatra
Description
The nominate race of red junglefowl has a mix of feather colours, with orange, brown, red, gold, grey, white, olive, and even metallic green plumage. The tail of the male rooster can grow up to , and the whole bird may be as long as . There are 14 tail feathers. A moult in June changes the bird's plumage to an eclipse pattern, which lasts through October. The male eclipse pattern includes a black feather in the middle of the back and small red-orange plumes spread across the body. Female eclipse plumage is generally indistinguishable from the plumage at other seasons, but the moulting schedule is the same as that of males.
Compared to the more familiar domestic chicken, the red junglefowl has a much smaller body mass (around lbs (1 kg) in females and lbs (1.5 kg) in males) and is brighter in coloration. Junglefowl are also behaviourally different from domestic chickens, being naturally very shy of humans compared to the much tamer domesticated subspecies.
Sexual dimorphism
Male junglefowl are significantly larger than females and have brightly coloured decorative feathers. The male's tail is composed of long, arching feathers that initially look black, but shimmer with blue, purple, and green in direct light. He also has long, golden hackle feathers on his neck and his back. The female's plumage is typical of this family of birds in being cryptic and adapted for camouflage. She alone looks after the eggs and chicks. She also has a very small comb and wattles (fleshy ornaments on the head that signal good health to rivals and potential mates) compared to the males.
During their mating season, the male birds announce their presence with the well-known "cock-a-doodle-doo" call or crowing. Within flocks, only dominant males crow. Male red junglefowl have a shorter crowing sound than domestic roosters; the call cuts off abruptly at the end. This serves both to attract potential mates and to make other male birds in the area aware of the risk of fighting a breeding competitor. A spur on the lower leg just behind and above the foot serves in such fighting. Their call structure is complex and they have distinctive alarm calls for aerial and ground predators to which others react appropriately.
Genetics
Orthology
G. gallus has three transferrins, all of which cluster closely with other vertebrates' orthologs.
Distribution and habitat
The range of the wild form stretches from India, Nepal and Bangladesh in the west, and eastwards across southern China, to Indochina; south/southeast into Malaysia, Singapore, the Philippines, and Indonesia. Junglefowl/Chickens were one of three main animals (along with domesticated pigs and dogs) carried by early Austronesian peoples from Island Southeast Asia in their voyages to the islands of Oceania in prehistory, starting around 5,000 years BP . Today, their modern descendants are found throughout Micronesia, Melanesia, and Polynesia.
Red junglefowl prefer disturbed habitats and edges, both natural and human-created. The forage and thick cover in these sorts of areas are attractive to junglefowl, especially nesting females. Junglefowl use both deforested and regenerating forests, and often are found near human settlements or areas of regrowth from slash-and-burn cultivation. Areas burned to promote bamboo growth also attract junglefowl, with edible bamboo seeds more available. In some areas, red junglefowl are absent from silvicultural and rubber plantations; elsewhere, they will occur in both tea plants and palm oil plantations. In the state of Selangor, Malaysia, palm foliage provides suitable cover; palm nut fruit provides adequate food, as well as insects (and their larvae) within, and adjacent to, the trees. The palms also offer an array of roost sites, from the low perches (~4 m) favored by females with chicks to the higher perches (up to 12 m) used by other adults.
Red junglefowl drink surface water when it is available, but they do not require it. Birds in North-Central India visit water holes frequently during the dry season, although not all junglefowl on the subcontinent live close enough to water to do so; population densities may thus be lower, where surface water is limited.
Behaviour and ecology
Red junglefowl regularly bathe in dust to keep the right balance of oil in their plumage. The dust absorbs extra oil and subsequently falls off.
Flight in these birds is almost purely confined to reaching their roosting areas at sunset in trees or any other high and relatively safe places free from ground predators, and for escape from immediate danger through the day.
Dominant male junglefowl appear to defend a territory against other dominant males, and the size of the territories has been inferred based on the proximity of roosts. Beebe concluded that territories were rather small, especially as compared to some of the pheasants with which he was familiar. This was supported by Collias and Collias, who reported that adjacent roost sites can be as close as . Within flocks, male red junglefowl exhibit dominance hierarchies, and dominant males tend to have larger combs than subordinate males. Red junglefowl typically live in flocks of one to a few males and several females. Males are more likely to occur alone than females.
Breeding
Males make a food-related display called "tidbitting", performed upon finding food in the presence of a female. The display is composed of coaxing, cluck-like calls, and eye-catching bobbing and twitching motions of the head and neck. During the performance, the male repeatedly picks up and drops the food item with his beak. The display usually ends when the hen takes the food item either from the ground or directly from the male's beak. Eventually, they sometimes mate.
In many areas, red junglefowl breed during the dry portion of the year, typically winter or spring. This is true in parts of India, Nepal, Thailand, Vietnam, and Laos. However, year-round breeding by red junglefowl has been documented in palm oil plantations in Malaysia and also may occur elsewhere. During the laying period, red junglefowl females lay an egg every day. Eggs take 21 days to develop. Chicks fledge in about 4 to 5 weeks, and at 12 weeks old they are chased out of the group by their mother — at which point they start a new group or join an existing one. Sexual maturity is reached at 5 months, with females taking slightly longer than males to reach maturity.
Dominant males attempt to maintain exclusive reproductive access to females, though females choose to mate with subordinate males about 40% of the time in a free-ranging feral flock in San Diego, California.
Diet
Red junglefowl are attracted to areas with ripe fruit or seeds, including fruit plantations, fields of domestic grain, and stands of bamboo. Although junglefowl typically eat fallen fruits and seeds on the ground, they occasionally forage in trees by perching on branches and feeding on hanging fruit. Fruits and seeds of scores of plant species have been identified from junglefowl crops, along with grasses, leaves, roots, and tubers. In addition, red junglefowl capture a wide variety of arthropods, other invertebrates, and vertebrates such as small lizards. Even mammalian faeces may be consumed. Many of these items are taken opportunistically as the birds forage, although some arthropods, such as termites, are taken in large quantities; about 1,000 individual termites have been found in a single crop. Plant materials constitute a higher proportion of the diet of adult red junglefowl than do arthropods and other animals. In contrast, chicks eat mostly adult and larval insects, earthworms, and only occasional plant material.
Relationship to humans
Chickens were created when red junglefowl were domesticated for human use around 8,000 years ago as subspecies Gallus gallus domesticus. They are now a major source of food for humans. However, undomesticated red junglefowl still represent an important source of meat and eggs in their endemic range. The undomesticated form is sometimes used in cock-fighting.
Timeline of domestication
In 2012, a study examined mitochondrial DNA recovered from ancient bones from Europe, Thailand, the Pacific, and Chile, and from Spanish colonial sites in Florida and the Dominican Republic, in directly dated samples originating in Europe at 1,000 BP and in the Pacific at 3,000 BP. The chicken was primarily domesticated from red junglefowl, with subsequent genetic contributions from grey junglefowl, Sri Lankan junglefowl, and green junglefowl. Domestication occurred about 8,000 years ago, as based on molecular evidence from a common ancestor flock in the bird's natural range, and then proceeded in waves both east and west. Zoogeography and evolutionary biology point to the original domestication site of chickens as somewhere in Mainland Southeast Asia and southern China in the Neolithic. Chickens were one of the ancestral domesticated animals of the Austronesian peoples. They were transported to Taiwan and the Philippines around 5,500 to 4,500 years ago. From there, they spread outwards with the Austronesian migrations to the rest of Island Southeast Asia, Micronesia, Island Melanesia, and Polynesia in prehistoric times.
Other archaeological evidence suggests domestication dates around 7,400 BP from the Chishan site, in the Hebei province of China. However, the domestication event in China has now been disputed by several studies citing unfavourable weather conditions at the time. In the Ganges region of India, wild red junglefowl were being used by humans as early as 7,000 years ago. No domestic chicken remains older than 4,000 years have been identified in the Indus Valley, and the antiquity of chickens recovered from excavations at Mohenjodaro is still debated.
Hybridization
The other three members of the genus — Sri Lanka junglefowl (G. lafayetii), grey junglefowl (G. sonneratii), and green junglefowl (G. varius) — do not usually produce fertile hybrids with the red junglefowl. However, supporting the hypothesis of a hybrid origin, research published in 2008 found that the gene responsible for the yellow skin of the domestic chicken most likely originated in the closely related grey junglefowl and not from the red junglefowl. Similarly, a 2020 study that analysed the whole genomes of Sri Lanka junglefowl, grey junglefowl, and green junglefowl found strong introgressive hybridisation events in different populations of indigenous village chickens. The study also shows that 71–79% of red junglefowl DNA is shared with the domestic chicken. A culturally significant hybrid between the red junglefowl and the green junglefowl in Indonesia is known as the bekisar.
Conservation status
Wild-type red junglefowl are thought to be facing threats due to hybridisation at forest edges, where domesticated free-ranging chickens are common. The red junglefowl is considered near threatened in Singapore. Nevertheless, they are classified by the IUCN as a species of least concern.
| Biology and health sciences | Galliformes | Animals |
470742 | https://en.wikipedia.org/wiki/Keratitis | Keratitis | Keratitis is a condition in which the eye's cornea, the clear dome on the front surface of the eye, becomes inflamed. The condition is often marked by moderate to intense pain and usually involves any of the following symptoms: pain, impaired eyesight, photophobia (light sensitivity), red eye and a 'gritty' sensation. Diagnosis of infectious keratitis is usually made clinically based on the signs and symptoms as well as eye examination, but corneal scrapings may be obtained and evaluated using microbiological culture or other testing to identify the causative pathogen.
Classification (by chronicity)
Acute
Acute epithelial keratitis
Nummular keratitis
Interstitial keratitis
Disciform keratitis
Chronic
Neurotrophic keratitis
Mucous plaque keratitis
Classification (infective)
Viral
The most common causes of viral keratitis include herpes simplex virus (HSV) and varicella zoster virus (VZV), which cause herpes simplex keratitis and herpes zoster keratitis (a subtype of herpes zoster ophthalmicus) respectively. Herpes simplex keratitis occurs due to latent HSV reactivation in the ophthalmic nerve (the V1 branch of the trigeminal nerve). Herpes keratitis is classically associated with a branching (dendritic) infiltrate pattern of inflammation in the corneal epithelium and may cause clouding of the cornea.
Approximately 8-20% of cases of shingles (due to VZV reactivation) involve the eyes as herpes zoster ophthalmicus. And VZV keratitis occurs in 13-76% of cases of herpes zoster ophthalmicus, usually 1 month after onset of symptoms. Herpes zoster ophthalmicus is also associated with reactivation of ZVZ in the V1 branch (the ophthalmic nerve) of the trigeminal nerve. VZV keratitis presents as a dendriform epithelial keratitis pattern early in the course of the infection. ZVZ keratitis may cause clouding of the cornea, with 50% of cases involving inflammatory infiltrates in the stromal layer of the cornea, corneal scarring is a possible complication of VZV keratitis. Vaccination with the zoster vaccine is highly effective in preventing shingles, as well as herpes zoster ophthalmicus and herpes zoster keratitis.
Bacterial
Bacterial keratitis. Bacterial infection of the cornea can follow from an injury or from wearing contact lenses. The bacteria involved are Staphylococcus aureus and for contact lens wearers, Pseudomonas aeruginosa. Pseudomonas aeruginosa produces enzymes that can digest the cornea. In those who wear contact lenses, bacteria are the most common causative agent of keratitis, with 90% of cases being due to a bacterial pathogen. Of those 90% of cases, Pseudomonas aeruginosa is responsible for 40%. Staph aureus and streptococci are other common bacterial pathogens responsible for infectious keratitis in contact lens wearers. Lens cases, used to store contact lenses, may form a biofilm leading to colonization of the contact lenses by bacteria, this is especially common with poor contact lens hygiene or improper storage.
Fungal
Fungal keratitis, caused by Aspergillus fumigatus and Candida albicans (cf. Fusarium, causing an outbreak of keratitis in 2005–2006 through the possible vector of Bausch & Lomb ReNu with MoistureLoc contact lens solution)
Amoebic
Acanthamoebic keratitis
Amoebic infection of the cornea is a serious corneal infection, most often affecting contact lens wearers. It is usually caused by Acanthamoeba. On May 25, 2007, the U.S. Center for Disease Control issued a health advisory due to increased risk of Acanthamoeba keratitis associated with use of Advanced Medical Optics Complete Moisture Plus Multi-Purpose eye solution.
Parasitic
Onchocercal keratitis, which follows Onchocerca volvulus infection by infected blackfly bite. These blackfly, Simulium'', usually dwell near fast-flowing African streams, so the disease is also called "river blindness".
Microbial keratitis (due to bacterial, fungal, or parasitic pathogens), as opposed to viral keratitis, is more commonly associated with the formation of corneal ulcers. Other risk factors for corneal ulcer formation include contact lens use, keratitis in the setting of eye trauma, underlying corneal disease or ocular surface diseases (such as severe chronic dry eye). Infectious keratitis sometimes presents as corneal edema, or with a hypopyon (a collection of inflammatory cells in the anterior chamber of the eye).
Classification (by stage of disease)
Superficial punctate keratitis
Ulcerative keratitis
Classification (by environmental aetiology)
Exposure keratitis (also known as exposure keratopathy) — due to dryness of the cornea caused by incomplete or inadequate eyelid closure (lagophthalmos).
Photokeratitis — keratitis due to intense ultraviolet radiation exposure (e.g. snow blindness or welder's arc eye.)
Contact lens acute red eye (CLARE) — a non-ulcerative sterile keratitis associated with colonization of Gram-negative bacteria on contact lenses.
Treatment
Treatment depends on the cause of the keratitis. Infectious keratitis can progress rapidly, and generally requires urgent antibacterial, antifungal, or antiviral therapy to eliminate the pathogen. Antibacterial solutions include levofloxacin, gatifloxacin, moxifloxacin, ofloxacin. It is unclear if steroid eye drops are useful.
In addition, contact lens wearers are typically advised to discontinue contact lens wear and replace contaminated contact lenses and contact lens cases. (Contaminated lenses and cases should not be discarded as cultures from these can be used to identify the pathogen).
Topical ganciclovir or oral valacyclovir, famciclovir or acyclovir are used for HSV keratitis. Steroids should be avoided as application of steroids to a dendritic ulcer caused by HSV may result in rapid and significant worsening of the ulcer to form an 'amoeboid' or 'geographic' ulcer, so named because of the ulcer's map like shape.
Prevention
In those who wear contact lenses, good lens hygiene and storage practices reduce the risk of keratitis. Specific lens care practices which may lead to infectious keratitis include wearing contact lenses overnight or in the shower, not replacing contact lens cases, storing lenses in tap water rather than contact lens solution and topping off lens solution rather than replacing it regularly. Improper lens storage may lead to bacterial biofilm formation in the contact lens case and subsequent colonization of the lenses by bacteria. Exposure of the lens to tap water through improper storage or use may lead to acanthamoeba infection, as the amoeba is commonly found in tap water.
Acyclovir prophylaxis has been found to reduce the risk of additional episodes of herpes simplex viral eye diseases (as well as oral or facial herpes) including a 50% reduction in the incidence of HSV keratitis. There was no rebound effect, or increased rate of HSV related eye disease upon stopping acyclovir prophylaxis.
Prognosis
Some infections may scar the cornea, thereby limiting vision. Others may result in perforation of the cornea, endophthalmitis (an infection inside the eye), or even loss of the eye. With proper medical attention, infections can usually be successfully treated without long-term visual loss.
Acanthamoebic and fungal keratitis are difficult to treat and are associated with a poor prognosis.
In non-humans
Feline eosinophilic keratitis — affecting cats and horses; possibly initiated by feline herpesvirus 1 or other viral infection.
| Biology and health sciences | Infectious diseases by site | Health |
470752 | https://en.wikipedia.org/wiki/Expectation%E2%80%93maximization%20algorithm | Expectation–maximization algorithm | In statistics, an expectation–maximization (EM) algorithm is an iterative method to find (local) maximum likelihood or maximum a posteriori (MAP) estimates of parameters in statistical models, where the model depends on unobserved latent variables. The EM iteration alternates between performing an expectation (E) step, which creates a function for the expectation of the log-likelihood evaluated using the current estimate for the parameters, and a maximization (M) step, which computes parameters maximizing the expected log-likelihood found on the E step. These parameter-estimates are then used to determine the distribution of the latent variables in the next E step. It can be used, for example, to estimate a mixture of gaussians, or to solve the multiple linear regression problem.
History
The EM algorithm was explained and given its name in a classic 1977 paper by Arthur Dempster, Nan Laird, and Donald Rubin. They pointed out that the method had been "proposed many times in special circumstances" by earlier authors. One of the earliest is the gene-counting method for estimating allele frequencies by Cedric Smith. Another was proposed by H.O. Hartley in 1958, and Hartley and Hocking in 1977, from which many of the ideas in the Dempster–Laird–Rubin paper originated. Another one by S.K Ng, Thriyambakam Krishnan and G.J McLachlan in 1977. Hartley’s ideas can be broadened to any grouped discrete distribution. A very detailed treatment of the EM method for exponential families was published by Rolf Sundberg in his thesis and several papers, following his collaboration with Per Martin-Löf and Anders Martin-Löf. The Dempster–Laird–Rubin paper in 1977 generalized the method and sketched a convergence analysis for a wider class of problems. The Dempster–Laird–Rubin paper established the EM method as an important tool of statistical analysis. | Mathematics | Statistics | null |
470843 | https://en.wikipedia.org/wiki/Fight-or-flight%20response | Fight-or-flight response | The fight-or-flight or the fight-flight-freeze-or-fawn (also called hyperarousal or the acute stress response) is a physiological reaction that occurs in response to a perceived harmful event, attack, or threat to survival. It was first described by Walter Bradford Cannon in 1915. His theory states that animals react to threats with a general discharge of the sympathetic nervous system, preparing the animal for fighting or fleeing. More specifically, the adrenal medulla produces a hormonal cascade that results in the secretion of catecholamines, especially norepinephrine and epinephrine. The hormones estrogen, testosterone, and cortisol, as well as the neurotransmitters dopamine and serotonin, also affect how organisms react to stress. The hormone osteocalcin might also play a part.
This response is recognised as the first stage of the general adaptation syndrome that regulates stress responses among vertebrates and other organisms.
Name
Originally understood as the "fight-or-flight" response in Cannon's research, the state of hyperarousal results in several responses beyond fighting or fleeing. This has led people to calling it the "fight, flight, freeze" response, "fight-flight-freeze-fawn"
or "fight-flight-faint-or-freeze", among other variants.
The wider array of responses, such as freezing, flop, faint, flee and fright, has led researchers to use more neutral or accommodating terminology such as "hyperarousal" or the "acute stress response".
Physiology
Autonomic nervous system
The autonomic nervous system is a control system that acts largely unconsciously and regulates heart rate, digestion, respiratory rate, pupillary response, urination, and sexual arousal. This system is the primary mechanism in control of the fight-or-flight response and its role is mediated by two different components: the sympathetic nervous system and the parasympathetic nervous system.
Sympathetic nervous system
The sympathetic nervous system originates in the spinal cord and its main function is to activate the arousal responses that occur during the fight-or-flight response. The sympathetic nervous system transfers signals from the dorsal hypothalamus, which activates the heart, increases vascular resistance, and increases blood flow, especially to the muscle, heart, and brain tissues. It activates the adrenal medulla, releasing catecholamines that amplify the sympathetic response. Additionally, this component of the autonomic nervous system utilizes and activates the release of norepinephrine by the adrenal glands in the reaction.
Parasympathetic nervous system
The parasympathetic nervous system originates in the sacral spinal cord and medulla, physically surrounding the sympathetic origin, and works in concert with the sympathetic nervous system. It is known as the calming portion of the autonomic nervous system. While the sympathetic nervous system is activated, the parasympathetic nervous system decreases its response. Efferent vagal fibers originating from the nucleus ambiguous fire in parallel to the respiratory system, decreasing the vagal cardiac parasympathetic tone. After the fight or flight response, the parasympathetic system's main function is to activate the "rest and digest" response and return the body to homeostasis. This system utilizes and activates the release of the neurotransmitter acetylcholine.
Reaction
The reaction begins in the amygdala, which triggers a neural response in the hypothalamus. The initial reaction is followed by activation of the pituitary gland and secretion of the hormone ACTH. The adrenal gland is activated almost simultaneously, via the sympathetic nervous system, and releases the hormone epinephrine. The release of chemical messengers results in the production of the hormone cortisol, which increases blood pressure, blood sugar, and suppresses the immune system.
The initial response and subsequent reactions are triggered in an effort to create a boost of energy. This boost of energy is activated by epinephrine binding to liver cells and the subsequent production of glucose. Additionally, the circulation of cortisol functions to turn fatty acids into available energy, which prepares muscles throughout the body for response.
Catecholamine hormones, such as adrenaline (epinephrine) or noradrenaline (norepinephrine), facilitate immediate physical reactions associated with a preparation for violent muscular action.
Function of physiological changes
The physiological changes that occur during the fight or flight response are activated to give the body increased strength and speed in anticipation of fighting or running. Some of the specific physiological changes and their functions include:
Increased blood flow to the muscles activated by diverting blood flow from other parts of the body to make taking quick action easier.
Increased blood pressure and heart rate enhance cardiac output in order to supply the body with more energy.
The liver secretes increased amounts of glucose (through adrenaline-induced glycogenolysis) and fats into the blood to provide the body with a fuel source to meet energy demands.
The respiratory rate increases to supply the oxygen necessary to help burn the extra glucose.
The blood clotting function of the body speeds up in order to reduce bleeding and prevent excessive blood loss in the event of an injury sustained during the response.
Increased muscle tension in order to provide the body with extra speed and strength, which can result in trembling or shaking until the tension is released.
The pupils dilate to let in more light, allowing for better vision of the body's surroundings.
Emotional components
Emotion regulation
In the context of the fight or flight response, emotional regulation is used proactively to avoid threats of stress or to control the level of emotional arousal. Emotional socialization can develop someone's ability to successfully regulate their emotions. Faced with a perceived threat (in the context of a fight or flight situation) those raised with supportive parental behaviors are far more likely to easily self-regulate their emotions.
Emotional reactivity
During the reaction, the intensity of emotion that is brought on by the stimulus will also determine the nature and intensity of the behavioral response. In a experiment conducted by Clayton, Lang, Leshner and Quick (2019), they viewed the responses of 49 participants to antitobacco messages. Participants reacted in two orders of fashion after seeing the message with the individual smoker and their affects on those surrounding them. The first reaction was participants who had higher defense mechanisms, who decided to ignore the messages, while the other participants who had lower defense mechanisms, ended up arguing and becoming frustrated after viewing the antitobacco messages. Individuals with higher levels of emotional reactivity (Such as an anxiety disorder) may be prone to anxiety and aggression, which illustrates the implications of appropriate emotional reaction in the fight or flight response.
Cognitive components
Content specificity
The specific components of cognitions in the fight or flight response seem to be largely negative. These negative cognitions may be characterised by: attention to negative stimuli, the perception of ambiguous situations as negative, and the recurrence of recalling negative words. There also may be specific negative thoughts associated with emotions commonly seen in the reaction.
Perception of control
Perceived control relates to an individual's thoughts about control over situations and events. Perceived control should be differentiated from actual control because an individual's beliefs about their abilities may not reflect their actual abilities. Therefore, overestimation or underestimation of perceived control can lead to anxiety and aggression.
Social information processing
The social information processing model proposes a variety of factors that determine behavior in the context of social situations and preexisting thoughts. The attribution of hostility, especially in ambiguous situations, seems to be one of the most important cognitive factors associated with the fight or flight response because of its implications towards aggression.
Other animals
Evolutionary perspective
An evolutionary psychology explanation is that early animals had to react to threatening stimuli quickly and did not have time to psychologically and physically prepare themselves. The fight or flight response provided them with the mechanisms to rapidly respond to threats against survival.
Examples
A typical example of the stress response is a grazing zebra. If the zebra sees a lion closing in for the kill, the stress response is activated as a means to escape its predator. The escape requires intense muscular effort, supported by all of the body's systems. The sympathetic nervous system's activation provides for these needs. A similar example involving fight is of a cat about to be attacked by a dog. The cat shows accelerated heartbeat, piloerection (hair standing on end), and pupil dilation, all signs of sympathetic arousal. Note that the zebra and cat still maintain homeostasis in all states.
In July 1992, Behavioral Ecology published experimental research conducted by biologist Lee A. Dugatkin where guppies were sorted into "bold", "ordinary", and "timid" groups based upon their reactions when confronted by a smallmouth bass (i.e. inspecting the predator, hiding, or swimming away) after which the guppies were left in a tank with the bass. After 60 hours, 40 percent of the timid guppies and 15 percent of the ordinary guppies survived while none of the bold guppies did.
Varieties of responses
Animals respond to threats in many complex ways. Rats, for instance, try to escape when threatened but will fight when cornered. Some animals stand perfectly still so that predators will not see them. Many animals freeze or play dead when touched in the hope that the predator will lose interest.
Other animals have alternative self-protection methods. Some species of cold-blooded animals change color swiftly to camouflage themselves. These responses are triggered by the sympathetic nervous system, but, in order to fit the model of fight or flight, the idea of flight must be broadened to include escaping capture either in a physical or sensory way. Thus, flight can be disappearing to another location or just disappearing in place, and fight and flight are often combined in a given situation.
The fight or flight actions also have polarity – the individual can either fight against or flee from something that is threatening, such as a hungry lion, or fight for or fly towards something that is needed, such as the safety of the shore from a raging river.
A threat from another animal does not always result in immediate fight or flight. There may be a period of heightened awareness, during which each animal interprets behavioral signals from the other. Signs such as paling, piloerection, immobility, sounds, and body language communicate the status and intentions of each animal. There may be a sort of negotiation, after which fight or flight may ensue, but which might also result in playing, mating, or nothing at all. An example of this is kittens playing: each kitten shows the signs of sympathetic arousal, but they never inflict real damage.
In criminal law
Acute stress response is a common issue in self-defense criminal cases. Expert opinions are usually required if the defender's fault becomes the focus of the case.
| Biology and health sciences | Ethology | Biology |
471376 | https://en.wikipedia.org/wiki/Pruning | Pruning | Pruning is a horticultural, arboricultural, and silvicultural practice involving the selective removal of certain parts of a plant, such as branches, buds, or roots.
The practice entails the targeted removal of diseased, damaged, dead, non-productive, structurally unsound, or otherwise unwanted plant material from crop and landscape plants. In general, the smaller the branch that is cut, the easier it is for a woody plant to compartmentalize the wound and thus limit the potential for pathogen intrusion and decay. It is therefore preferable to make any necessary formative structural pruning cuts to young plants, rather than removing large, poorly placed branches from mature plants.
Woody plants may undergo a process referred to as "self-pruning", where they will drop twigs or branches which are no longer producing more energy than they require. It is theorized that this process can also occur in response to lack of water, in order to reduce the surface area where water can be lost. This natural shedding of branches is called cladoptosis.
Specialized pruning practices may be applied to certain plants, such as roses, fruit trees, and grapevines. Different pruning techniques may be used on herbaceous plants than those used on perennial woody plants.
Reasons to prune plants include deadwood removal, shaping (by controlling or redirecting growth), improving or sustaining health, reducing risk from falling branches, preparing nursery specimens for transplanting, and both harvesting and increasing the yield or quality of flowers and fruits.
Pruning terms
Branch wood
Branch wood is an individual stem that grows off of another stem.
Trunk wood
Trunk wood is the main stem of a tree which individual stems grow out of.
Branch collar
This refers to the area below the union of where branch wood attaches with the trunk/stem wood. This can often appear raised.
Branch bark ridge
This refers to the junction between branch wood and trunk/stem wood. It usually looks raised.
Types of Pruning
Pruning in an urban setting is crucial due to the tree being in drastically different conditions than where they naturally grow.
Arborists, orchardists, and gardeners use various garden tools and tree cutting tools designed for the purpose, such as secateurs, loppers, handsaws, or chainsaws. Additionally in forestry, pole pruners and pole saws are commonly used and these are often attached to poles that reach up to 5-6 m, this is a more efficient way of pruning than with ladders. These bush saws on polls have also been motorized as chainsaws which is even more efficient. Older technology used Billhooks, Kaiser blades and pruning knives. Although still used in some coppicing they are not used so much in commercial forestry due to the difficulty of cutting flush with the stem. Flush cuts happen when you cut into the cambium layer of the main trunk which can happen when you are not precise with pruning cuts and remove a portion of the branch collar which can put the tree at risk of entry cords from forest pathogens.
Although there are several different types of pruning they can be simplified into two categories. One of which is cutting the branch back to a specific and intermediate point, called reduction cut, and the other completely removes a branch back to the union where the branch connects which the main trunk, called removal cut.
Reduction cuts is when you remove a portion of a growing stem down to a set of desirable buds or side-branching stems. This is commonly performed in well trained plants for a variety of reasons, for example to stimulate growth of flowers, fruit or branches, as a preventive measure to wind and snow damage on long stems and branches, and finally to encourage growth of the stems in a desirable direction.
Thinning: A more drastic form of pruning, a thinning out cut is the removal of an entire shoot, limb, or branch at its point of origin. This is usually employed to revitalize a plant by removing over-mature, weak, problematic, and excessive growths. When performed correctly, thinning encourages the formation of new growth that will more readily bear fruit and flowers. This is a common technique in pruning roses and for amplifying and "opening-up" the branching of neglected trees, or for renewing shrubs with multiple branches.
Topping: Topping is a very severe form of pruning which involves removing all branches and growths down to a few large branches or to the trunk of the tree. When performed correctly it is used on very young trees, and can be used to begin training younger trees for pollarding or for trellising to form an espalier.
Raising removes the lower branches from a tree in order to provide clearance for buildings, vehicles, pedestrians, and vistas.
Reduction reduces the size of a tree, often for clearance for utility lines. Reducing the height or spread of a tree is best accomplished by pruning back the leaders and branch terminals to lateral branches that are large enough to assume the terminal roles (at least one-third the diameter of the cut stem). Compared to topping, reduction helps maintain the form and structural integrity of the tree.
In orchards, fruit trees are often lopped to encourage regrowth and to maintain a smaller tree for ease of picking fruit. The pruning regime in orchards is more planned and the productivity of each tree is an important factor.
Deadwooding
Branches die off for a number of reasons including sunlight deficiency, pest and disease damage, and root structure damage. A dead branch will at some point decay back to the parent stem and fall off. This is normally a slow process but can be hastened by high winds or extreme temperatures. The main reason deadwooding is performed is safety. Situations that usually demand removal of deadwood include trees that overhang public roads, houses, public areas, power lines, telephone cables and gardens. Trees located in wooded areas are usually assessed as lower risk but assessments consider the number of visitors. Trees adjacent to footpaths and access roads are often considered for deadwood removal.
Another reason for deadwooding is amenity value, i.e. a tree with a large amount of deadwood throughout the crown will look more aesthetically pleasing with the deadwood removed. The physical practice of deadwooding can be carried out most of the year though should be avoided when the tree is coming into leaf. The deadwooding process speeds up the tree's natural abscission process. It also reduces unwanted weight and wind resistance and can help overall balance.
Preventive structural pruning
Preventative and structural pruning can be done to mitigate several issues young trees may have in the future. The structural pruning can reduce tree stress, increase the lifespan of trees, and promotes resistance to damage due to natural weather events. Attributes of trees with good structure include excurrent growth by having a single dominant leader, branch unions without included bark, and a balanced canopy. Structural pruning does this by developing or maintaining a dominant leader, identify the lowest branches in the canopy, prevent branches below the permanent canopy from growing too large, keeping all branches less than one half the trunk diameter, space main branches along one dominant trunk, and suppress growth on branches with included bark.
Subordination pruning
Subordination pruning is done on limbs that will exceed 50% percent of the stem diameter. A reduction cut may be performed while still allowing about 50% of the branch. This is done to help maintain form and deter the formation of co-dominant leaders. Temporary branches may be too large for a removal cut so subordination pruning should be done to slowly reduce a limb by 50% each year to allow the tree to properly heal from the cut. As a tree becomes larger the slower it grows. Reducing the larger limbs for eventual removal will allow for the tree to promote new growth rather than using energy in encouraging unwanted limbs to continue to grow. Removing a large branch increases the likelihood of the cut to not heal properly which also may attract insects, diseases and fungus.
Crown thinning
Crown thinning is the removal of live healthy branches which increases light penetration, air circulation and reduces wind resistance which reduces risks from damage and the possibility of pest infestation.
Crown raising
Crown raising involves the removal of the lower branches to a given height. The height is achieved by the removal of whole branches or removing the parts of branches which extend below the desired height. The branches are normally not lifted to more than one third of the tree's total height.
Crown lifting is done for access; these being pedestrian, vehicle or space for buildings and street furniture. Lifting the crown will allow traffic and pedestrians to pass underneath safely. This pruning technique is usually used in the urban environment as it is for public safety and aesthetics rather than tree form and timber value.
Crown lifting introduces light to the lower part of the trunk; this, in some species can encourage epicormic growth from dormant buds. To reduce this sometimes smaller branches are left on the lower part of the trunk. Excessive removal of the lower branches can displace the canopy weight, this will make the tree top heavy, therefore adding stress to the tree. When a branch is removed from the trunk, it creates a large wound. This wound is susceptible to disease and decay, and could lead to reduced trunk stability. Therefore, much time and consideration must be taken when choosing the height the crown is to be lifted to.
This would be an inappropriate operation if the tree species’ form was of a shrubby nature. This would therefore remove most of the foliage and would also largely unbalance the tree. This procedure should not be carried out if the tree is in decline, poor health or dead, dying or dangerous (DDD) as the operation will remove some of the photosynthetic area the tree uses. This will increase the decline rate of the tree and could lead to death.
If the tree is of great importance to an area or town, (i.e. veteran or ancient) then an alternative solution to crown lifting would be to move the target or object so it is not in range. For example, diverting a footpath around a tree's drip line so the crown lift is not needed. Another solution would be to prop up or cable-brace the low hanging branch. This is a non-invasive solution which in some situations may be more economical and environmentally friendly.
Vista pruning
Selectively pruning a window of view in a tree.
Crown reduction
Reducing the height and or spread of a tree by selectively cutting back to smaller branches and in fruit trees for increasing of light interception and enhancing fruit quality.
Pollarding
A regular form of pruning where certain deciduous species are pruned back to pollard heads every year in the dormant period. This practice is usually commenced on juvenile trees so they can adapt to the harshness of the practice. This practice can be used for tree shaping but is also used in specific species which young branches can be sold for floral arrangements.
Deadheading
Deadheading is the act of removing spent flowers or flowerheads for aesthetics, to prolong bloom for up to several weeks or promote rebloom, or to prevent seeding.
Time period
In general, pruning deadwood and small branches can be done at any time of year. Depending on the species, many temperate plants can be pruned either during dormancy in winter, or, for species where winter frost can harm a recently pruned plant, after flowering is completed. In the temperate areas of the northern hemisphere autumn pruning should be avoided, as the spores of disease and decay fungi are abundant at this time of year.
Some woody plants tend to bleed profusely from cuts, such as mesquite and maple. Some callus over slowly, such as magnolia. In this case, they are better pruned during active growth when they can more readily heal. Woody plants that flower early in the season, on spurs that form on wood that has matured the year before, such as apples, should be pruned right after flowering as later pruning will sacrifice flowers the following season. Forsythia, azaleas and lilacs all fall into this category.
| Technology | Horticulture | null |
471385 | https://en.wikipedia.org/wiki/Impatiens | Impatiens | Impatiens is a genus of more than 1,000 species of flowering plants, widely distributed throughout the Northern Hemisphere and the tropics. Together with the genus Hydrocera (one species), Impatiens make up the family Balsaminaceae.
Common names in North America include impatiens, jewelweed, touch-me-not, snapweed and patience. As a rule-of-thumb, "jewelweed" is used exclusively for Nearctic species, and balsam is usually applied to tropical species. In the British Isles by far the most common names are impatiens and busy lizzie, especially for the many varieties, hybrids and cultivars involving Impatiens walleriana. "Busy lizzie" is also found in the American literature. Impatiens glandulifera is commonly called policeman's helmet in the UK, where it is an introduced species.
Description
Most Impatiens species are herbaceous annuals or perennials with succulent stems. Only a few woody species exist. Plant size varies, from five centimeters to 2.5 meters, depending on the species. Stems often form roots when they come into contact with the soil. The leaves are entire, often dentate or sinuate with extrafloral nectaries. Depending on the species, leaves can be thin to succulent. Particularly on the underside of the leaves, tiny air bubbles are trapped over and under the leaf surface, giving them a silvery sheen that becomes pronounced when they are held underwater.
The zygomorphic flowers of Impatiens are protandric (male becoming female with age). The calyx consists of five free sepals, of which one pair is often strongly reduced. The non-paired sepal forms a flower spur-producing nectar. In a group of species from Madagascar, the spur is completely lacking, but they still have three sepals. The crown consists of five petals, of which the lateral pairs are fused. The five stamens are fused and form a cap over the ovary, which falls off after the male phase. After the stamens have fallen off, the female phase starts and the stigma becomes receptive, which reduces self-pollination.
The scientific name Impatiens (Latin for "impatient") and the common name "touch-me-not" refer to the explosive dehiscence of the seed capsules. The mature capsules burst, sending seeds up to several meters away.
Distribution
The genus Impatiens occurs in Africa, Eurasia and North America. Two species (Impatiens turrialbana and Impatiens mexicana) occur in isolated areas in Central America (southern Mexico and Costa Rica). Most Impatiens species occur in the tropical and subtropical mountain forests in Africa, Madagascar, the Himalayas, the Western Ghats (southwest India) and southeast Asia. In Europe only a single Impatiens species (Impatiens noli-tangere) occurs naturally. However, several neophytic species exist.
In the 19th and 20th centuries, humans transported the North American orange jewelweed (I. capensis) to England, France, the Netherlands, Poland, Sweden, Finland, and potentially other areas of Northern and Central Europe. For example, it was not recorded from Germany as recently as 1996, but since then a population has naturalized in Hagen at the Ennepe River. The orange jewelweed is quite similar to the touch-me-not balsam (I. noli-tangere), the only Impatiens species native to Central and Northern Europe, and it utilizes similar habitats, but no evidence exists of natural hybrids between them. Small balsam (I. parviflora), originally native to southern Central Asia, is even more extensively naturalized in Europe. More problematic is the Himalayan balsam (I. glandulifera), a densely growing species which displaces smaller plants by denying them sunlight. It is an invasive weed in many places, and tends to dominate riparian vegetation along polluted rivers and nitrogen-rich spots. Thus, it exacerbates ecosystem degradation by forming stands where few other plants can grow, and by rendering riverbanks more prone to erosion, as it has only a shallow root system.
Ecology
Most Impatiens species are perennial herbs. However, several annual species exist, especially in the temperate regions as well as in the Himalayas. A few Impatiens species in southeast Asia (e.g. Impatiens kerriae or Impatiens mirabilis) form shrubs or small trees up to three meters tall. Most Impatiens species occur in forests, especially along streams and paths or at the forest edge with a little bit of sunlight. Additionally, a few species occur in open landscapes, such as heathland, river banks or savanna.
The genus Impatiens is characterized by a large variety of flower architectures. Traditionally two flower types are differentiated: one with a sacculate spur and a more or less two-lipped flower and a second with a filiform spur and a flat flower surface. However, several transition forms exist. Additionally, a group of 125 spur-less species exist on Madagascar, forming a third main flower type. Due to the large variability in flower architecture it seems reasonable to group the species by their main pollinators: such as bees and bumblebees, butterflies, moths, flies, and sunbirds. Further, a few cleistogamous species exist. However, most species are dependent on pollinator activity for efficient seed production but many of them are self-compatible. Most temperate species as well as some tropical species can switch from chasmogamous (pollinator-dependent) to cleistogamous (seed production within closed flowers) flowers when nutrient and light conditions become adverse.
Impatiens foliage is used for food by the larvae of some Lepidoptera species, such as the dot moth (Melanchra persicariae), as well as other insects, such as the Japanese beetle (Popillia japonica). The leaves are toxic to many other animals, including the budgerigar (Melopsittacus undulatus), but the bird will readily eat the flowers. The flowers are also visited by bumblebees and certain Lepidoptera, such as the common spotted flat (Celaenorrhinus leucocera).
Parasitic plants that use impatiens as hosts include the European dodder (Cuscuta europaea). A number of plant diseases affect this genus.
The starkly differing flower shapes found in this genus, combined with the easy cultivation of many species, have served to make some balsam species model organisms in plant evolutionary developmental biology. Also, Impatiens is rather closely related to the carnivorous plant families Roridulaceae and Sarraceniaceae. Peculiar stalked glands found on balsam sepals secrete mucus and might be related to the structures from which the prey-catching and -digesting glands of these carnivorous plants evolved. Balsams are not known to be protocarnivorous plants, however.
In 2011–2013, the United States experienced a significant outbreak of the fungal disease downy mildew that affects impatiens, particularly Impatiens walleriana. The disease was also reported in Canada as well. The pathogen plasmopara obducens is the chief culprit suspected by scientists, but Bremiella sphaerosperma is related. These pathogens were first reported in the United States in 2004.
Medicinal uses and phytochemistry
Impatiens contain 2-methoxy-1,4-naphthoquinone, an anti-inflammatory and fungicide naphthoquinone that is an active ingredient in some formulations of Preparation H.
North American impatiens have been used as herbal remedies for the treatment of bee stings, insect bites, and stinging nettle (Urtica dioica) rashes. They are also used after poison ivy (Toxicodendron radicans) contact to prevent a rash from developing. The efficacy of orange jewelweed (I. capensis) and yellow jewelweed (I. pallida) in preventing poison ivy contact dermatitis has been studied, with conflicting results. A study in 1958 found that Impatiens biflora was an effective alternative to standard treatment for dermatitis caused by contact with sumac, while later studies found that the species had no antipruritic effects after the rash has developed. Researchers reviewing these contradictions state that potential reason for these conflicts include the method of preparation and timing of application. A 2012 study found that while an extract of orange jewelweed and garden jewelweed (I. balsamina) was not effective in reducing contact dermatitis, a mash of the plants applied topically decreased it.
Impatiens glandulifera is one of the Bach flower remedies, flower extracts used as herbal remedies for physical and emotional problems. It is included in the "Rescue Remedy" or "Five Flower Remedy", a potion touted as a treatment for acute anxiety and which is supposed to be protective in stressful situations. Studies have found no difference between the effect of the potion and that of a placebo.
All Impatiens taste bitter and seem to be slightly toxic upon ingestion, causing intestinal ailments like vomiting and diarrhea. The toxic compounds have not been identified but are probably the same as those responsible for the bitter taste, likely might be glycosides or alkaloids.
α-Parinaric acid, a polyunsaturated fatty acid discovered in the seeds of the makita tree (Atuna racemosa subsp. racemosa), is together with linolenic acid the predominant component of the seed fat of garden jewelweed (I. balsamina), and perhaps other species of Impatiens. This is interesting from a phylogenetic perspective, because the makita tree is a member of the Chrysobalanaceae in a lineage of eudicots entirely distinct from the balsams.
Certain jewelweeds, including the garden jewelweed contain the naphthoquinone lawsone, a dye that is also found in henna (Lawsonia inermis) and is also the hair coloring and skin coloring agent in mehndi. In ancient China, Impatiens petals mashed with rose and orchid petals and alum were used as nail polish: leaving the mixture on the nails for some hours colored them pink or reddish.
Cultivation
Impatiens are popular garden annuals. Hybrids, typically derived from busy lizzie (I. walleriana) and New Guinea impatiens (I. hawkeri), have commercial importance as garden plants. I. walleriana is native to East Africa, and yielded 'Elfin' series of cultivars, which was subsequently improved as the 'Super Elfin' series. Double-flowered cultivars also exist.
Other Impatiens species, such as African queen (I. auricoma), garden jewelweed (I. balsamina), blue diamond impatiens (I. namchabarwensis), parrot flower (I. psittacina), Congo cockatoo (I. niamniamensis), Ceylon balsam (I. repens), and poor man's rhododendron (I. sodenii), are also used as ornamental plants.
Species
| Biology and health sciences | Ericales | null |
471476 | https://en.wikipedia.org/wiki/Astronomical%20system%20of%20units | Astronomical system of units | The astronomical system of units, formerly called the IAU (1976) System of Astronomical Constants, is a system of measurement developed for use in astronomy. It was adopted by the International Astronomical Union (IAU) in 1976 via Resolution No. 1, and has been significantly updated in 1994 and 2009 (see Astronomical constant).
The system was developed because of the difficulties in measuring and expressing astronomical data in International System of Units (SI units). In particular, there is a huge quantity of very precise data relating to the positions of objects within the Solar System that cannot conveniently be expressed or processed in SI units. Through a number of modifications, the astronomical system of units now explicitly recognizes the consequences of general relativity, which is a necessary addition to the International System of Units in order to accurately treat astronomical data.
The astronomical system of units is a tridimensional system, in that it defines units of length, mass and time. The associated astronomical constants also fix the different frames of reference that are needed to report observations. The system is a conventional system, in that neither the unit of length nor the unit of mass are true physical constants, and there are at least three different measures of time.
Astronomical unit of time
The astronomical unit of time is the day, defined as seconds. 365.25 days make up one Julian year. The symbol D is used in astronomy to refer to this unit.
Astronomical unit of mass
The astronomical unit of mass is the solar mass. The symbol is often used to refer to this unit.
The solar mass (), , is a standard way to express mass in astronomy, used to describe the masses of other stars and galaxies. It is equal to the mass of the Sun, about times the mass of the Earth or 1 048 times the mass of Jupiter.
In practice, the masses of celestial bodies appear in the dynamics of the Solar System only through the products GM, where G is the constant of gravitation. In the past, GM of the Sun could be determined experimentally with only limited accuracy. Its present accepted value is .
Jupiter mass
Jupiter mass ( or MJUP), is the unit of mass equal to the total mass of the planet Jupiter, . Jupiter mass is used to describe masses of the gas giants, such as the outer planets and extrasolar planets. It is also used in describing brown dwarfs and Neptune-mass planets.
Earth mass
Earth mass () is the unit of mass equal to that of the Earth. 1 = . Earth mass is often used to describe masses of rocky terrestrial planets. It is also used to describe Neptune-mass planets. One Earth mass is times a Jupiter mass.
Astronomical unit of length
The astronomical unit of length is now defined as exactly 149 597 870 700 meters. It is approximately equal to the mean Earth–Sun distance. It was formerly defined as that length for which the Gaussian gravitational constant (k) takes the value when the units of measurement are the astronomical units of length, mass and time. The dimensions of k2 are those of the constant of gravitation (G), i.e., L3M−1T−2. The term "unit distance" is also used for the length A while, in general usage, it is usually referred to simply as the "astronomical unit", symbol au.
An equivalent formulation of the old definition of the astronomical unit is the radius of an unperturbed circular Newtonian orbit about the Sun of a particle having infinitesimal mass, moving with a mean motion of radians per day.
The speed of light in IAU is the defined value c0 = of the SI units. In terms of this speed, the old definition of the astronomical unit of length had the accepted value: 1 au = c0τA = () m, where τA is the transit time of light across the astronomical unit. The astronomical unit of length was determined by the condition that the measured data in the ephemeris match observations, and that in turn decides the transit time τA.
Other units for astronomical distances
The distances to distant galaxies are typically not quoted in distance units at all, but rather in terms of redshift. The reasons for this are that converting redshift to distance requires knowledge of the Hubble constant, which was not accurately measured until the early 21st century, and that at cosmological distances, the curvature of spacetime allows one to come up with multiple definitions for distance. For example, the distance as defined by the amount of time it takes for a light beam to travel to an observer is different from the distance as defined by the apparent size of an object.
| Physical sciences | Measurement systems | Basics and measurement |
471487 | https://en.wikipedia.org/wiki/Cosmic%20distance%20ladder | Cosmic distance ladder | The cosmic distance ladder (also known as the extragalactic distance scale) is the succession of methods by which astronomers determine the distances to celestial objects. A direct distance measurement of an astronomical object is possible only for those objects that are "close enough" (within about a thousand parsecs) to Earth. The techniques for determining distances to more distant objects are all based on various measured correlations between methods that work at close distances and methods that work at larger distances. Several methods rely on a standard candle, which is an astronomical object that has a known luminosity.
The ladder analogy arises because no single technique can measure distances at all ranges encountered in astronomy. Instead, one method can be used to measure nearby distances, a second can be used to measure nearby to intermediate distances, and so on. Each rung of the ladder provides information that can be used to determine the distances at the next higher rung.
Direct measurement
At the base of the ladder are fundamental distance measurements, in which distances are determined directly, with no physical assumptions about the nature of the object in question. The precise measurement of stellar positions is part of the discipline of astrometry. Early fundamental distances—such as the radii of the earth, moon and sun, and the distances between them—were well estimated with very low technology by the ancient Greeks.
Astronomical unit
Direct distance measurements are based upon the astronomical unit (AU), which is defined as the mean distance between the Earth and the Sun.
Kepler's laws provide precise ratios of the orbit sizes of objects orbiting the Sun, but provide no measurement of the overall scale of the orbit system. Radar is used to measure the distance between the orbits of the Earth and of a second body. From that measurement and the ratio of the two orbit sizes, the size of Earth's orbit is calculated. The Earth's orbit is known with an absolute precision of a few meters and a relative precision of a few parts in 100 billion ().
Historically, observations of Venus transits were crucial in determining the AU; in the first half of the 20th century, observations of asteroids were also important. Presently the orbit of Earth is determined with high precision using radar measurements of distances to Venus and other nearby planets and asteroids, and by tracking interplanetary spacecraft in their orbits around the Sun through the Solar System.
Parallax
Standard candles
Almost all astronomical objects used as physical distance indicators belong to a class that has a known brightness. By comparing this known luminosity to an object's observed brightness, the distance to the object can be computed using the inverse-square law. These objects of known brightness are termed standard candles, coined by Henrietta Swan Leavitt.
The brightness of an object can be expressed in terms of its absolute magnitude. This quantity is derived from the logarithm of its luminosity as seen from a distance of 10 parsecs. The apparent magnitude, the magnitude as seen by the observer (an instrument called a bolometer is used), can be measured and used with the absolute magnitude to calculate the distance d to the object in parsecs as follows:
or
where m is the apparent magnitude, and M the absolute magnitude. For this to be accurate, both magnitudes must be in the same frequency band and there can be no relative motion in the radial direction.
Some means of correcting for interstellar extinction, which also makes objects appear fainter and more red, is needed, especially if the object lies within a dusty or gaseous region. The difference between an object's absolute and apparent magnitudes is called its distance modulus, and astronomical distances, especially intergalactic ones, are sometimes tabulated in this way.
Problems
Two problems exist for any class of standard candle. The principal one is calibration, that is the determination of exactly what the absolute magnitude of the candle is. This includes defining the class well enough that members can be recognized, and finding enough members of that class with well-known distances to allow their true absolute magnitude to be determined with enough accuracy. The second problem lies in recognizing members of the class, and not mistakenly using a standard candle calibration on an object which does not belong to the class. At extreme distances, which is where one most wishes to use a distance indicator, this recognition problem can be quite serious.
A significant issue with standard candles is the recurring question of how standard they are. For example, all observations seem to indicate that Type Ia supernovae that are of known distance have the same brightness, corrected by the shape of the light curve. The basis for this closeness in brightness is discussed below; however, the possibility exists that the distant Type Ia supernovae have different properties than nearby Type Ia supernovae. The use of Type Ia supernovae is crucial in determining the correct cosmological model. If indeed the properties of Type Ia supernovae are different at large distances, i.e. if the extrapolation of their calibration to arbitrary distances is not valid, ignoring this variation can dangerously bias the reconstruction of the cosmological parameters, in particular the reconstruction of the matter density parameter.
That this is not merely a philosophical issue can be seen from the history of distance measurements using Cepheid variables. In the 1950s, Walter Baade discovered that the nearby Cepheid variables used to calibrate the standard candle were of a different type than the ones used to measure distances to nearby galaxies. The nearby Cepheid variables were population I stars with much higher metal content than the distant population II stars. As a result, the population II stars were actually much brighter than believed, and when corrected, this had the effect of doubling the estimates of distances to the globular clusters, the nearby galaxies, and the diameter of the Milky Way.
Most recently kilonova have been proposed as another type of standard candle. "Since kilonovae explosions are spherical, astronomers could compare the apparent size of a supernova explosion with its actual size as seen by the gas motion, and thus measure the rate of cosmic expansion at different distances."
Standard siren
Gravitational waves originating from the inspiral phase of compact binary systems, such as neutron stars or black holes, have the useful property that energy emitted as gravitational radiation comes exclusively from the orbital energy of the pair, and the resultant shrinking of their orbits is directly observable as an increase in the frequency of the emitted gravitational waves. To leading order, the rate of change of frequency is given by
where is the gravitational constant, is the speed of light, and is a single (therefore computable) number called the chirp mass of the system, a combination of the masses of the two objects
By observing the waveform, the chirp mass can be computed and thence the power (rate of energy emission) of the gravitational waves. Thus, such a gravitational wave source is a standard siren of known loudness.
Just as with standard candles, given the emitted and received amplitudes, the inverse-square law determines the distance to the source. There are some differences with standard candles, however. Gravitational waves are not emitted isotropically, but measuring the polarisation of the wave provides enough information to determine the angle of emission. Gravitational wave detectors also have anisotropic antenna patterns, so the position of the source on the sky relative to the detectors is needed to determine the angle of reception.
Generally, if a wave is detected by a network of three detectors at different locations, the network will measure enough information to make these corrections and obtain the distance. Also unlike standard candles, gravitational waves need no calibration against other distance measures. The measurement of distance does of course require the calibration of the gravitational wave detectors, but then the distance is fundamentally given as a multiple of the wavelength of the laser light being used in the gravitational wave interferometer.
There are other considerations that limit the accuracy of this distance, besides detector calibration. Fortunately, gravitational waves are not subject to extinction due to an intervening absorbing medium. But they are subject to gravitational lensing, in the same way as light. If a signal is strongly lensed, then it might be received as multiple events, separated in time, the analogue of multiple images of a quasar, for example. Less easy to discern and control for is the effect of weak lensing, where the signal's path through space is affected by many small magnification and demagnification events. This will be important for signals originating at cosmological redshifts greater than 1. It is difficult for detector networks to measure the polarization of a signal accurately if the binary system is observed nearly face-on. Such signals suffer significantly larger errors in the distance measurement. Unfortunately, binaries radiate most strongly perpendicular to the orbital plane, so face-on signals are intrinsically stronger and the most commonly observed.
If the binary consists of a pair of neutron stars, their merger will be accompanied by a kilonova/hypernova explosion that may allow the position to be accurately identified by electromagnetic telescopes. In such cases, the redshift of the host galaxy allows a determination of the Hubble constant . This was the case for GW170817, which was used to make the first such measurement. Even if no electromagnetic counterpart can be identified for an ensemble of signals, it is possible to use a statistical method to infer the value of .
Standard ruler
Another class of physical distance indicator is the standard ruler. In 2008, galaxy diameters have been proposed as a possible standard ruler for cosmological parameter determination. More recently the physical scale imprinted by baryon acoustic oscillations (BAO) in the early universe has been used.
In the early universe (before recombination) the baryons and photons scatter off each other, and form a tightly coupled fluid that can support sound waves. The waves are sourced by primordial density perturbations, and travel at speed that can be predicted from the baryon density and other cosmological parameters.
The total distance that these sound waves can travel before recombination determines a fixed scale, which simply expands with the universe after recombination. BAO therefore provide a standard ruler that can be measured in galaxy surveys from the effect of baryons on the clustering of galaxies. The method requires an extensive galaxy survey in order to make this scale visible, but has been measured with percent-level precision (see baryon acoustic oscillations). The scale does depend on cosmological parameters like the baryon and matter densities, and the number of neutrinos, so distances based on BAO are more dependent on cosmological model than those based on local measurements.
Light echos can be also used as standard rulers, although it is challenging to correctly measure the source geometry.
Galactic distance indicators
With few exceptions, distances based on direct measurements are available only out to about a thousand parsecs, which is a modest portion of our own Galaxy. For distances beyond that, measures depend upon physical assumptions, that is, the assertion that one recognizes the object in question, and the class of objects is homogeneous enough that its members can be used for meaningful estimation of distance.
Physical distance indicators, used on progressively larger distance scales, include:
Dynamical parallax, uses orbital parameters of visual binaries to measure the mass of the system, and hence use the mass–luminosity relation to determine the luminosity
Eclipsing binaries — In the last decade, measurement of eclipsing binaries' fundamental parameters has become possible with 8-meter class telescopes. This makes it feasible to use them as indicators of distance. Recently, they have been used to give direct distance estimates to the Large Magellanic Cloud (LMC), Small Magellanic Cloud (SMC), Andromeda Galaxy and Triangulum Galaxy. Eclipsing binaries offer a direct method to gauge the distance to galaxies to a new improved 5% level of accuracy which is feasible with current technology to a distance of around 3 Mpc (3 million parsecs).
RR Lyrae variables — used for measuring distances within the galaxy and in nearby globular clusters.
The following four indicators all use stars in the old stellar populations (Population II):
Tip of the red-giant branch (TRGB) distance indicator.
Planetary nebula luminosity function (PNLF)
Globular cluster luminosity function (GCLF)
Surface brightness fluctuation (SBF)
In galactic astronomy, X-ray bursts (thermonuclear flashes on the surface of a neutron star) are used as standard candles. Observations of X-ray burst sometimes show X-ray spectra indicating radius expansion. Therefore, the X-ray flux at the peak of the burst should correspond to Eddington luminosity, which can be calculated once the mass of the neutron star is known (1.5 solar masses is a commonly used assumption). This method allows distance determination of some low-mass X-ray binaries. Low-mass X-ray binaries are very faint in the optical, making their distances extremely difficult to determine.
Interstellar masers can be used to derive distances to galactic and some extragalactic objects that have maser emission.
Cepheids and novae
The Tully–Fisher relation
The Faber–Jackson relation
Type Ia supernovae that have a very well-determined maximum absolute magnitude as a function of the shape of their light curve and are useful in determining extragalactic distances up to a few hundred Mpc. A notable exception is SN 2003fg, the "Champagne Supernova", a Type Ia supernova of unusual nature.
Redshifts and Hubble's law
Main sequence fitting
When the absolute magnitude for a group of stars is plotted against the spectral classification of the star, in a Hertzsprung–Russell diagram, evolutionary patterns are found that relate to the mass, age and composition of the star. In particular, during their hydrogen burning period, stars lie along a curve in the diagram called the main sequence. By measuring these properties from a star's spectrum, the position of a main sequence star on the H–R diagram can be determined, and thereby the star's absolute magnitude estimated. A comparison of this value with the apparent magnitude allows the approximate distance to be determined, after correcting for interstellar extinction of the luminosity because of gas and dust.
In a gravitationally-bound star cluster such as the Hyades, the stars formed at approximately the same age and lie at the same distance. This allows relatively accurate main sequence fitting, providing both age and distance determination.
Extragalactic distance scale
The extragalactic distance scale is a series of techniques used today by astronomers to determine the distance of cosmological bodies beyond our own galaxy, which are not easily obtained with traditional methods. Some procedures use properties of these objects, such as stars, globular clusters, nebulae, and galaxies as a whole. Other methods are based more on the statistics and probabilities of things such as entire galaxy clusters.
Wilson–Bappu effect
Discovered in 1956 by Olin Wilson and M.K. Vainu Bappu, the Wilson–Bappu effect uses the effect known as spectroscopic parallax. Many stars have features in their spectra, such as the calcium K-line, that indicate their absolute magnitude. The distance to the star can then be calculated from its apparent magnitude using the distance modulus.
There are major limitations to this method for finding stellar distances. The calibration of the spectral line strengths has limited accuracy and it requires a correction for interstellar extinction. Though in theory this method has the ability to provide reliable distance calculations to stars up to 7 megaparsecs (Mpc), it is generally only used for stars at hundreds of kiloparsecs (kpc).
Classical Cepheids
Beyond the reach of the Wilson–Bappu effect, the next method relies on the period-luminosity relation of classical Cepheid variable stars. The following relation can be used to calculate the distance to Galactic and extragalactic classical Cepheids:
Several problems complicate the use of Cepheids as standard candles and are actively debated, chief among them are: the nature and linearity of the period-luminosity relation in various passbands and the impact of metallicity on both the zero-point and slope of those relations, and the effects of photometric contamination (blending) and a changing (typically unknown) extinction law on Cepheid distances.
These unresolved matters have resulted in cited values for the Hubble constant ranging between 60 km/s/Mpc and 80 km/s/Mpc. Resolving this discrepancy is one of the foremost problems in astronomy since some cosmological parameters of the Universe may be constrained significantly better by supplying a precise value of the Hubble constant.
Cepheid variable stars were the key instrument in Edwin Hubble's 1923 conclusion that M31 (Andromeda) was an external galaxy, as opposed to a smaller nebula within the Milky Way. He was able to calculate the distance of M31 to 285 kpc, today's value being 770 kpc.
As detected thus far, NGC 3370, a spiral galaxy in the constellation Leo, contains the farthest Cepheids yet found at a distance of 29 Mpc. Cepheid variable stars are in no way perfect distance markers: at nearby galaxies they have an error of about 7% and up to a 15% error for the most distant.
Supernovae
There are several different methods for which supernovae can be used to measure extragalactic distances.
Measuring a supernova's photosphere
We can assume that a supernova expands in a spherically symmetric manner. If the supernova is close enough such that we can measure the angular extent, θ(t), of its photosphere, we can use the equation
where ω is angular velocity, θ is angular extent. In order to get an accurate measurement, it is necessary to make two observations separated by time Δt. Subsequently, we can use
where d is the distance to the supernova, Vej is the supernova's ejecta's radial velocity (it can be assumed that Vej equals Vθ if spherically symmetric).
This method works only if the supernova is close enough to be able to measure accurately the photosphere. Similarly, the expanding shell of gas is in fact not perfectly spherical nor a perfect blackbody. Also interstellar extinction can hinder the accurate measurements of the photosphere. This problem is further exacerbated by core-collapse supernova. All of these factors contribute to the distance error of up to 25%.
Type Ia light curves
Type Ia supernovae are some of the best ways to determine extragalactic distances. Ia's occur when a binary white dwarf star begins to accrete matter from its companion star. As the white dwarf gains matter, eventually it reaches its Chandrasekhar limit of .
Once reached, the star becomes unstable and undergoes a runaway nuclear fusion reaction. Because all Type Ia supernovae explode at about the same mass, their absolute magnitudes are all the same. This makes them very useful as standard candles. All Type Ia supernovae have a standard blue and visual magnitude of
Therefore, when observing a Type Ia supernova, if it is possible to determine what its peak magnitude was, then its distance can be calculated. It is not intrinsically necessary to capture the supernova directly at its peak magnitude; using the multicolor light curve shape method (MLCS), the shape of the light curve (taken at any reasonable time after the initial explosion) is compared to a family of parameterized curves that will determine the absolute magnitude at the maximum brightness. This method also takes into effect interstellar extinction/dimming from dust and gas.
Similarly, the stretch method fits the particular supernovae magnitude light curves to a template light curve. This template, as opposed to being several light curves at different wavelengths (MLCS) is just a single light curve that has been stretched (or compressed) in time. By using this Stretch Factor, the peak magnitude can be determined.
Using Type Ia supernovae is one of the most accurate methods, particularly since supernova explosions can be visible at great distances (their luminosities rival that of the galaxy in which they are situated), much farther than Cepheid Variables (500 times farther). Much time has been devoted to the refining of this method. The current uncertainty approaches a mere 5%, corresponding to an uncertainty of just 0.1 magnitudes.
Novae in distance determinations
Novae can be used in much the same way as supernovae to derive extragalactic distances. There is a direct relation between a nova's max magnitude and the time for its visible light to decline by two magnitudes. This relation is shown to be:
Where is the time derivative of the nova's mag, describing the average rate of decline over the first 2 magnitudes.
After novae fade, they are about as bright as the most luminous Cepheid variable stars, therefore both these techniques have about the same max distance: ~ 20 Mpc. The error in this method produces an uncertainty in magnitude of about ±0.4
Globular cluster luminosity function
Based on the method of comparing the luminosities of globular clusters (located in galactic halos) from distant galaxies to that of the Virgo Cluster, the globular cluster luminosity function carries an uncertainty of distance of about 20% (or 0.4 magnitudes).
US astronomer William Alvin Baum first attempted to use globular clusters to measure distant elliptical galaxies. He compared the brightest globular clusters in Virgo A galaxy with those in Andromeda, assuming the luminosities of the clusters were the same in both. Knowing the distance to Andromeda, Baum has assumed a direct correlation and estimated Virgo A's distance.
Baum used just a single globular cluster, but individual formations are often poor standard candles. Canadian astronomer René Racine assumed the use of the globular cluster luminosity function (GCLF) would lead to a better approximation. The number of globular clusters as a function of magnitude is given by:
where m0 is the turnover magnitude, M0 is the magnitude of the Virgo cluster, and sigma is the dispersion ~ 1.4 mag.
It is assumed that globular clusters all have roughly the same luminosities within the universe. There is no universal globular cluster luminosity function that applies to all galaxies.
Planetary nebula luminosity function
Like the GCLF method, a similar numerical analysis can be used for planetary nebulae within far off galaxies. The planetary nebula luminosity function (PNLF) was first proposed in the late 1970s by Holland Cole and David Jenner. They suggested that all planetary nebulae might all have similar maximum intrinsic brightness, now calculated to be M = −4.53. This would therefore make them potential standard candles for determining extragalactic distances.
Astronomer George Howard Jacoby and his colleagues later proposed that the PNLF function equaled:
Where N(M) is number of planetary nebula, having absolute magnitude M. M* is equal to the nebula with the brightest magnitude.
Surface brightness fluctuation method
The following method deals with the overall inherent properties of galaxies. These methods, though with varying error percentages, have the ability to make distance estimates beyond 100 Mpc, though it is usually applied more locally.
The surface brightness fluctuation (SBF) method takes advantage of the use of CCD cameras on telescopes. Because of spatial fluctuations in a galaxy's surface brightness, some pixels on these cameras will pick up more stars than others. As distance increases, the picture will become increasingly smoother. Analysis of this describes a magnitude of the pixel-to-pixel variation, which is directly related to a galaxy's distance.
Sigma-D relation
The Sigma-D relation (or Σ-D relation), used in elliptical galaxies, relates the angular diameter (D) of the galaxy to its velocity dispersion. It is important to describe exactly what D represents, in order to understand this method. It is, more precisely, the galaxy's angular diameter out to the surface brightness level of 20.75 B-mag arcsec−2. This surface brightness is independent of the galaxy's actual distance from us. Instead, D is inversely proportional to the galaxy's distance, represented as d. Thus, this relation does not employ standard candles. Rather, D provides a standard ruler. This relation between D and Σ is
where C is a constant which depends on the distance to the galaxy clusters.
This method has the potential to become one of the strongest methods of galactic distance calculators, perhaps exceeding the range of even the Tully–Fisher method. As of today, however, elliptical galaxies are not bright enough to provide a calibration for this method through the use of techniques such as Cepheids. Instead, calibration is done using more crude methods.
Overlap and scaling
A succession of distance indicators, which is the distance ladder, is needed for determining distances to other galaxies. The reason is that objects bright enough to be recognized and measured at such distances are so rare that few or none are present nearby, so there are too few examples close enough with reliable trigonometric parallax to calibrate the indicator. For example, Cepheid variables, one of the best indicators for nearby spiral galaxies, cannot yet be satisfactorily calibrated by parallax alone, though the Gaia space mission can now weigh in on that specific problem. The situation is further complicated by the fact that different stellar populations generally do not have all types of stars in them.
Cepheids in particular are massive stars, with short lifetimes, so they will only be found in places where stars have very recently been formed. Consequently, because elliptical galaxies usually have long ceased to have large-scale star formation, they will not have Cepheids. Instead, distance indicators whose origins are in an older stellar population (like novae and RR Lyrae variables) must be used. RR Lyrae variables are less luminous than Cepheids, and novae are unpredictable and an intensive monitoring program—and luck during that program—is needed to gather enough novae in the target galaxy for a good distance estimate.
Because the more distant steps of the cosmic distance ladder depend upon the nearer ones, the more distant steps include the effects of errors in the nearer steps, both systematic and statistical ones. The result of these propagating errors means that distances in astronomy are rarely known to the same level of precision as measurements in the other sciences, and that the precision necessarily is poorer for more distant types of object.
Another concern, especially for the very brightest standard candles, is their "standardness": how homogeneous the objects are in their true absolute magnitude. For some of these different standard candles, the homogeneity is based on theories about the formation and evolution of stars and galaxies, and is thus also subject to uncertainties in those aspects. For the most luminous of distance indicators, the Type Ia supernovae, this homogeneity is known to be poor. However, no other class of object is bright enough to be detected at such large distances, so the class is useful simply because there is no real alternative.
The observational result of Hubble's law, the proportional relationship between distance and the speed with which a galaxy is moving away from us, usually referred to as redshift, is a product of the cosmic distance ladder. Edwin Hubble observed that fainter galaxies are more redshifted. Finding the value of the Hubble constant was the result of decades of work by many astronomers, both in amassing the measurements of galaxy redshifts and in calibrating the steps of the distance ladder. Hubble's law is the primary means we have for estimating the distances of quasars and distant galaxies in which individual distance indicators cannot be seen.
| Physical sciences | Astrometry | null |
12840258 | https://en.wikipedia.org/wiki/Nuclear%20envelope | Nuclear envelope | The nuclear envelope, also known as the nuclear membrane, is made up of two lipid bilayer membranes that in eukaryotic cells surround the nucleus, which encloses the genetic material.
The nuclear envelope consists of two lipid bilayer membranes: an inner nuclear membrane and an outer nuclear membrane. The space between the membranes is called the perinuclear space. It is usually about 10–50 nm wide. The outer nuclear membrane is continuous with the endoplasmic reticulum membrane. The nuclear envelope has many nuclear pores that allow materials to move between the cytosol and the nucleus. Intermediate filament proteins called lamins form a structure called the nuclear lamina on the inner aspect of the inner nuclear membrane and give structural support to the nucleus.
Structure
The nuclear envelope is made up of two lipid bilayer membranes, an inner nuclear membrane and an outer nuclear membrane. These membranes are connected to each other by nuclear pores. Two sets of intermediate filaments provide support for the nuclear envelope. An internal network forms the nuclear lamina on the inner nuclear membrane. A looser network forms outside to give external support. The actual shape of the nuclear envelope is irregular. It has invaginations and protrusions and can be observed with an electron microscope.
Outer membrane
The outer nuclear membrane also shares a common border with the endoplasmic reticulum. While it is physically linked, the outer nuclear membrane contains proteins found in far higher concentrations than the endoplasmic reticulum. All four nesprin proteins (nuclear envelope spectrin repeat proteins) present in mammals are expressed in the outer nuclear membrane. Nesprin proteins connect cytoskeletal filaments to the nucleoskeleton. Nesprin-mediated connections to the cytoskeleton contribute to nuclear positioning and to the cell’s mechanosensory function. KASH domain proteins of Nesprin-1 and -2 are part of a LINC complex (linker of nucleoskeleton and cytoskeleton) and can bind directly to cystoskeletal components, such as actin filaments, or can bind to proteins in the perinuclear space. Nesprin-3 and -4 may play a role in unloading enormous cargo; Nesprin-3 proteins bind plectin and link the nuclear envelope to cytoplasmic intermediate filaments. Nesprin-4 proteins bind the plus end directed motor kinesin-1. The outer nuclear membrane is also involved in development, as it fuses with the inner nuclear membrane to form nuclear pores.
Inner membrane
The inner nuclear membrane encloses the nucleoplasm, and is covered by the nuclear lamina, a mesh of intermediate filaments which stabilizes the nuclear membrane as well as being involved in chromatin function. It is connected to the outer membrane by nuclear pores which penetrate the membranes. While the two membranes and the endoplasmic reticulum are linked, proteins embedded in the membranes tend to stay put rather than dispersing across the continuum. It is lined with a fiber network called the nuclear lamina which is 10-40 nm thick and provides strength.
Mutations in the genes that encode for the inner nuclear membrane proteins can cause several laminopathies.
Nuclear pores
The nuclear envelope is punctured by around a thousand nuclear pore complexes, about 100 nm across, with an inner channel about 40 nm wide. The complexes contain a number of nucleoporins, proteins that link the inner and outer nuclear membranes.
Cell division
During the G2 phase of interphase, the nuclear membrane increases its surface area and doubles its number of nuclear pore complexes.
In eukaryotes such as yeast which undergo closed mitosis, the nuclear membrane stays intact during cell division. The spindle fibers either form within the membrane, or penetrate it without tearing it apart.
In other eukaryotes (animals as well as plants), the nuclear membrane must break down during the prometaphase stage of mitosis to allow the mitotic spindle fibers to access the chromosomes inside. The breakdown and reformation processes are not well understood.
Breakdown
In mammals, the nuclear membrane can break down within minutes, following a set of steps during the early stages of mitosis.
First, M-Cdk's phosphorylate nucleoporin polypeptides and they are selectively removed from the nuclear pore complexes. After that, the rest of the nuclear pore complexes break apart simultaneously. Biochemical evidence suggests that the nuclear pore complexes disassemble into stable pieces rather than disintegrating into small polypeptide fragments. M-Cdk's also phosphorylate elements of the nuclear lamina (the framework that supports the envelope) leading to the disassembly of the lamina and hence the envelope membranes into small vesicles.
Electron and fluorescence microscopy has given strong evidence that the nuclear membrane is absorbed by the endoplasmic reticulum—nuclear proteins not normally found in the endoplasmic reticulum show up during mitosis.
In addition to the breakdown of the nuclear membrane during the prometaphase stage of mitosis, the nuclear membrane also ruptures in migrating mammalian cells during the interphase stage of the cell cycle. This transient rupture is likely caused by nuclear deformation. The rupture is rapidly repaired by a process dependent on "endosomal sorting complexes required for transport" (ESCRT) made up of cytosolic protein complexes. During nuclear membrane rupture events, DNA double-strand breaks occur. Thus the survival of cells migrating through confined environments appears to depend on efficient nuclear envelope and DNA repair machineries.
Aberrant nuclear envelope breakdown has also been observed in laminopathies and in cancer cells leading to mislocalization of cellular proteins, the formation of micronuclei and genomic instability.
Reformation
Exactly how the nuclear membrane reforms during telophase of mitosis is debated. Two theories exist—
Vesicle fusion — where vesicles of nuclear membrane fuse together to rebuild the nuclear membrane
Re-shaping of the endoplasmic reticulum—where the parts of the endoplasmic reticulum containing the absorbed nuclear membrane envelop the nuclear space, reforming a closed membrane.
Origin of the nuclear membrane
A study of the comparative genomics, evolution and origins of the nuclear membrane led to the proposal that the nucleus emerged in the primitive eukaryotic ancestor (the “prekaryote”), and was triggered by the archaeo-bacterial symbiosis. Several ideas have been proposed for the evolutionary origin of the nuclear membrane. These ideas include the invagination of the plasma membrane in a prokaryote ancestor, or the formation of a genuine new membrane system following the establishment of proto-mitochondria in the archaeal host. The adaptive function of the nuclear membrane may have been to serve as a barrier to protect the genome from reactive oxygen species (ROS) produced by the cells' pre-mitochondria.
| Biology and health sciences | Organelles | Biology |
12841199 | https://en.wikipedia.org/wiki/Baton%20%28law%20enforcement%29 | Baton (law enforcement) | A baton (also truncheon, nightstick, billy club, billystick, cosh, lathi, or simply stick) is a roughly cylindrical club made of wood, rubber, plastic, or metal. It is carried as a compliance tool and defensive weapon by law-enforcement officers, correctional staff, security guards and military personnel. The name baton comes from the French bâton (stick), derived from Old French Baston, from Latin bastum.
As a weapon a baton may be used defensively (to block) or offensively (to strike, jab, or bludgeon), and it can aid in the application of armlocks. The usual striking or bludgeoning action is not produced by a simple and direct hit, as with an ordinary blunt object, but rather by bringing the arm down sharply while allowing the truncheon to pivot nearly freely forward and downward, so moving its tip much faster than its handle. Batons are also used for non-weapon purposes such as breaking windows to free individuals trapped in a vehicle, or turning out a suspect's pockets during a search (as a precaution against sharp objects).
Some people other than law enforcement officers use batons as weapons because of their simple construction and easy concealment. The use or carrying of batons or improvised clubs by people other than law enforcement officers is restricted by law in many countries.
History
In the Victorian era, police in London carried truncheons about one foot long called billy clubs. According to the Online Etymology Dictionary, this name was first recorded in 1848 as slang for a burglars' crowbar. The meaning "policeman's club" is first recorded 1856. The truncheon acted as the policeman's 'Warrant Card' as the Royal Crest attached to it indicated the policeman's authority. This was always removed when the equipment left official service (often with the person who used it). Earlier on, the word was used in vulgar Latin (bastο—a stick helping walking, from basta—hold).
The Victorian original has since developed into the multitude of varieties available today. The typical truncheon is a straight stick made from wood or a synthetic material, approximately in diameter and long, with a fluted handle to aid in gripping. Truncheons are often ornamented with their organizations' coats of arms. Longer truncheons are called "riot batons" because of their use in riot control. Truncheons may have developed as a marriage between the club or military mace and the staff of office/sceptre.
Straight batons of rubber have a softer impact. Some of the kinetic energy bends and compresses the rubber and bounces off when the object is struck. Rubber batons are not very effective when used on the subject's arms or legs, and can still cause injury if the head is struck. That is why most police departments have stopped issuing them. The Russian police standard-issue baton is rubber, except in places such as Siberia, where it can be cold enough that the rubber may become brittle and break if struck. The traffic baton is red to make it more visible as a signaling aid in directing traffic. In Russia, traffic batons are striped in black and white for the same reason, and in Sweden they are white. Until the mid-1990s, British police officers carried traditional wooden truncheons of a sort that had changed little from Victorian times. Since the late 1990s, the collapsible baton is issued except for public order duties, where a fixed, acrylic baton is used. Side-handled batons were issued for a while, but fell out of favour.
The New York City Police Department used to use two kinds of batons depending on the time. The one for daytime was called a day-stick and was in length. Another baton, that was used at night, was long and called a night-stick, which is the origin of the word nightstick. The night-stick was longer so it could provide extra protection which was thought to be necessary at night.
Target areas
In modern police training, the primary targets are large nerve clusters, such as the common peroneal nerve in the mid-thigh and large, easily targetable muscle groups, such as the quadriceps and biceps. The baton is swung in fast, "snapping" strikes to these areas, sometimes only making contact with the tip. Taken together, these are intended to impair the subject's ability to continue advancing (by striking the leg) or attack (by striking the arm) by causing transitory neurapraxia (temporary muscle pain, spasm and paralysis due to nerve injury). Modern systems strictly prohibit hitting the skull, sternum, spine, or groin unless such an attack is conducted in defense of life, with many jurisdictions considering this deadly force.
Before the 1970s, a common use of the police baton was to strike a suspect's head with a full-force overhand motion in order to stun them or knock them unconscious by cerebral concussion, similar to the pre-baton practice of buffaloing with the handle of a revolver. However, this practice had two major liabilities. First, there was a high risk and incidence of death or permanent injury, as the difference in force between that required to concuss a suspect into non-resistance and that which would fracture their skull tends to be narrow and unpredictable. Second, there were problems with reliability, as resistance to cerebral concussion varies widely between individuals, and head strikes that did not disable the suspect were found to merely escalate the encounter. Officer Arthur Lamb, a well-known trainer on the baton, once stated:
As a result, civil lawsuits and claims of police brutality resulted in revised training for officers.
Designs
Batons in common use by police around the world include many different designs, such as fixed-length straight batons, blackjacks, fixed-length side-handle batons, collapsible straight batons, and other more exotic variations. All types have their advantages and disadvantages. The design and popularity of specific types of baton have evolved over the years and are influenced by a variety of factors. These include inherent compromises in the dual (and competing) goals of control effectiveness and safety (for both officer and subject).
Straightstick
A straight, fixed-length baton (also commonly referred to as a "straightstick") is the oldest and simplest police baton design, known as far back as ancient Egypt. It consists of little more than a long cylinder with a molded, turned or wrapped grip, usually with a slightly thicker or tapering shaft and rounded tip. They are often made of hardwood, but in modern times are available in other materials such as aluminium, acrylic, and dense plastics and rubber. They range in size from short clubs less than in length to "riot batons" commonly used in civil disturbances or by officers mounted on horseback.
Straightsticks tend to be heavier and have more weight concentrated in the striking end than other designs. This makes them less maneuverable, but theoretically would deliver more kinetic energy on impact. Most agencies have replaced the straightstick with other batons because of inconvenience to carry, and a desire for their officers to look less threatening to the community they serve. Despite having been replaced by side-handle and expandable batons in many (if not most) law enforcement agencies, straightsticks remain in use by many major departments in the US, such as the Baltimore, Denver, Sacramento, Long Beach, Santa Ana, Philadelphia, San Francisco, and Riverside Police Departments. They also are used by NYPD Auxiliary Police officers, as well as many Military Police forces around the world.
Side-handle
Side-handle batons (sometimes referred to as T-batons) are batons with a short side handle at a right angle to the shaft, about from one end. The main shaft is typically in length. They are derived from the tonfa, an Okinawan kobudō weapon, and are used with a similar technique (although tonfa are usually used in pairs, whereas side-handle batons are not). The best-known example is the Monadnock PR-24; "PR-24" has become a genericized trademark within the law enforcement and security communities for this type of product.
It can be held by:
One end, and the intersection between the shaft and the handle used to catch a long swung blunt or sharp weapon.
The side handle, and the long shaft held against the hand and forearm to splint and shield the arm against an expected blow from an attacker.
Side-handle batons are made in both fixed and collapsible models and may be constructed from a range of materials including wood, poly-carbonate, epoxy, aluminium, or a combination of materials. Some side-handle batons are one-piece design; the side-handle component and primary shaft are permanently fused together during manufacturing. One-piece designs are potentially stronger than two-piece designs, and have no risk of having a locking screw loosen from its threads. Other side-handle batons are two-piece in design (common among cheaper makes); the side-handle component is screwed into the primary shaft. The side handle may be removed from the shaft by the end-user, converting the side-handle into a straight baton.
The advantages of a side-handle baton over a straight baton are numerous:
There is a far greater number of defensive techniques/maneuvers that may be used with the side-handle baton in contrast with the straight baton.
The side-handle component may aid in weapon retention, making it more difficult for a suspect to take the baton away from the officer in a struggle.
The side-handle component prevents the baton from rolling far away if inadvertently dropped, unlike a straight baton.
Subjectively, some officers may be able to deliver a strike of greater power with the side-handle baton (when used in conjunction with a "power stroke") over a straight baton.
Due to its design, a side handle baton is generally used in a more defensive and less offensive manner than a straight baton, and thus it is less likely for an officer to "instinctively" use a side-handle baton as a simple bludgeon and direct indiscriminate strikes against a suspect. Also, the typically defensive stance the side-handle baton is used with is generally believed to present a more community-friendly image than a straight baton.
Side-handle batons have a few disadvantages:
More training is required for an officer to fully utilize the potential of a side-handle baton compared to a straight baton. Its use is not as instinctive as the straight batons.
The side-handle slightly increases overall weight and bulk of the baton compared to a straight baton of identical length.
When the side-handle baton is used as a simple bludgeon (without gripping the side-handle), it is less effective than a straight baton. Also, if an officer uses it incorrectly and strikes the subject with the end of the side handle while holding the baton from the other end it can cause serious injury (especially if the end of the handle is a large ball).
Side-handle batons have been involved in high-profile incidents of alleged police brutality, such as in New Zealand's 1981 Springbok Tour and the Rodney King beating.
Rapid Rotation Baton
It is a version of the side-handle baton released in the mid-1990s. It tries to address some of the disadvantages of straight, side and expandable batons and combine them with the strengths of the aforementioned.
Expandable
An expandable baton (also referred to variously as a collapsible baton, telescopic baton, tactical baton, spring cosh, ASP, or extendable) is typically composed of a cylindrical outer shaft containing telescoping inner shafts (typically 2 or 3, depending on the design) that lock into each other when expanded. The shafts are usually made of steel, but lightweight baton models may have their shafts made from other materials such as aluminum alloy.
Expandable batons may have a solid tip at the outer end of the innermost shaft; the purpose of the solid tip is to maximize the power of a strike when the baton is used as an impact weapon. Expandable batons are made in both straight and side-handle configurations but are considerably more common in the straight configuration. The best-known example of the straight expandable baton is the ASP Baton, from Armament Systems and Procedures.
Depending on the holster or scabbard design, it may be possible to carry an expandable baton in either collapsed or expanded position, which would be helpful if an officer needed to holster an expanded baton and it was not possible or convenient to collapse it at the time. An expandable baton is opened by being swung forcefully while collapsed, using inertia to extend and lock the segments by friction. Some mechanical-lock versions can also be opened by simply pulling the segments apart. Depending on the design, expandable batons may be collapsed either by being brought down (inverted) on a hard surface or by depressing a button lock and manually collapsing the shafts. Additionally, the baton, in the collapsed configuration, may be used as a control device against non-compliant subjects in conjunction with pain-compliance control techniques, such as to remove a driver refusing to exit his or her vehicle. It can be used as a large kubotan.
The expandable baton is provided to most officers in the British police forces, the idea being that should violence suddenly escalate, the baton can be easily deployed but can be stowed neatly away so as not to affect movement due to its mounting point on the officer's clothing. It is also commonly used in the UK and many other countries as a means of gaining entry quickly to a vehicle that contains offenders. In such a situation the baton is deployed and, due to the solid end of the device, is used to strike windows and/or windscreens of the vehicle to either gain entry or to stop the driver from seeing where they are going in circumstances where the officer has hit the screen while the vehicle is still in motion.
Advantages
The advantages of a collapsible baton over a fixed baton are numerous:
The collapsible shaft makes it easier for the officer to carry it and to sit in a car seat wearing it since when collapsed it is between long. This is contrasted with non-collapsible batons, which the officer may, as a measure of convenience, often resort to removing from their belt when seating themselves in a vehicle. This often results in leaving the baton behind when an officer is exiting the vehicle, and not readily expecting trouble.
Non-collapsible batons are typically carried in a ring-type belt attachment. Fixed batons carried in such holders may easily fall out of the holder when the officer wearing the baton sprints. Neither holding the baton down in the ring with a hand nor holding the baton in the hand is a good solution. The typical collapsible straight baton and its scabbard do not suffer this and remain secure regardless of the wearing officer's movement.
In theory, the mere display of extending the baton may in some instances be terrifying to an aggressive person (due to both the sight and sound of the action, with a similar intimidation technique as used with pump-action shotguns), and may thus de-escalate the situation through fear-motivated submission of the target without physical violence.
A collapsible baton may be deployed against a suspect whether expanded or collapsed; expanded, the baton's reach is extended, but collapsed, the baton is handier in close quarters. This provides greater versatility in a wider range of environments over the fixed-length baton.
A collapsible baton is essentially a heavy steel rod with usually a slightly wider tip, that concentrates the force of a blow more effectively and to a smaller area than a polycarbonate baton. This results in a strike that impacts harder to the muscle and causes deeper pain, removing the need for several strikes when targeting large muscle groups. Striking bones will cause serious damage.
Disadvantages
Expandable batons have some disadvantages:
Some police may prefer to carry a fixed baton due to the greater intimidation it may provide. Similarly, a fixed baton serves better as a conspicuous symbol of authority (i.e., "badge of office") than a collapsed expandable baton.
Fixed batons may often be less expensive than their collapsible counterparts of identical or similar quality. Because of this, some law enforcement departments, such as the Los Angeles County Sheriff's Department, may issue a fixed-length baton but have their officers/deputies purchase expandable batons at the option and expense of the individual officer.
Fixed batons may be inherently faster to bring into action because they do not need to be extended before usage as an impact weapon (unless one wishes to use a collapsible baton in collapsed form). It is, however, possible to deliver a strike whilst opening the baton in one fluid motion if the officer is correctly trained. This is called a "rapid response strike".
If an expandable baton is of friction-lock design, as most are, there is an inherent risk that the baton may inadvertently close at an inopportune moment while being used to strike. This also prevents expandable batons from being used to prod or thrust.
In a situation in which stealth is required, a collapsed baton may rattle, revealing the officer's position.
Most expandable batons have most of their weight concentrated at the grip and the tip tends to be the lightest part since it is the thinnest part of the baton. As such it may deliver less forceful blows than a fixed baton.
Multi-tool batons
Since early law enforcement professionals were very reliant on their batons there was a popular movement to outfit police batons with implements like whistles, torches/flashlights and tear gas. At least four models were built with weapon-retention devices that would deploy "sharp spikes or blades" in case a suspect tried to grab an officer's baton. There is even a mention of two early 20th century patents for batons with guns built into the frame in American Police Equipment.
Similar weapons
Blackjacks and saps
The terms blackjack, cosh, and sap refer to any of several short, easily concealed club weapons consisting of a dense (often lead) weight attached to the end of a short shaft. A form of bludgeon, these use a handle to accelerate the dense core and transfer kinetic energy from the swing to it. The blow can be directed at any normal blunt force target, but delivered to the head induces concussion. While usually meant to stun or knock out the subject, head strikes have a high risk of causing a permanent, disabling brain injury or a fatality.
The terminology applied to these weapons can be imprecise, depends on the source and time period, and may refer to an improvised device accomplishing the same result. Blackjacks have long been a favorite of the criminal element as a concealed carry weapon, and are illegal to privately own in many jurisdictions.
A late 19th-century type is a wooden shaft about one foot long, with a leather- or macramé-covered lead ball as the head. This weapon is referred to by some sources as a "sap" (derived from "sapling" due to its wood handle), or euphemistically as a "life-preserver". The term "cosh" may also originate with this weapon, being derived from the Romani word kašt, meaning "stick" or "piece of wood". The term "blackjack" referring to a hand weapon is of unknown etymology, and the earliest text reference is 1889.
A type used by 19th and early 20th century sailors for both self-defense and aggression was weighted with a lead ball at one or both ends of a piece of baleen, which was then wrapped in woven or plaited marline or codline and varnished over. Some carefully made examples were likely to have been used by a boatswain or ship's master-at-arms or ship's mate as a badge of office and discipline-enforcer, so some modern sources call this weapon a "bosun's cosh". The term "blackjack" is sometimes applied by early 20th-century maritime sources to a lead weight knotted or woven into the end of a short piece of rope that serves as a handle, though most sources would consider this weapon a type of slungshot.
In the 20th century newer shorter (and more readily concealable) designs emerged that were predominantly made of stitched or braided leather, with a flexible spring inside the handle to impart a whip-like action to compensate for leverage and acceleration sacrificed to reduced overall length. To balance forces and minimize unintended damage, this style came to employ a flexible material such as lead shot as a weight.
Law enforcement sources from the mid-20th century preferred to divide these into two categories: "blackjacks", which have a mostly cylindrical striking head, and "saps" which have a flat, usually oval-shaped head. In common usage, however, these terms have become interchangeable, so a "sap" of this latter kind is sometimes more precisely called a flat sap, slap jack or beavertail sap to differentiate it. The sap's flat profile makes it easier to carry in a pocket and spreads its impact out over a broader area, making it less likely to break bone. However, it can also be used to strike with the edge for more focused impact, though this was discouraged by most police departments for precisely this reason. Alternatively, some variants use lead shot, powdered metal, or even sand for the weight inside the head, usually called a "soft sap", which reduces the likelihood of bone fractures, particularly of the skull.
Blackjacks and saps were popular among law enforcement for a time due to their low profile, small size, and effectiveness in very close range, such as when grappling with a suspect. Besides the head, they were also used on the elbows, wrists, shins, collarbone, and groin. The flat sap, in particular, could be used to strike large muscle groups with the edge. In the early days of use, they were favored for their ability to stun or knock a suspect unconscious with a blow to the head. By the late 1960s head-strikes with impact weapons in general were strongly discouraged by most police departments and trainers because of the risk of death or permanent injury, as well as questionable effectiveness. By the 1990s virtually all modern police departments had phased them out from their issued equipment, and most banned their use entirely.
Stunguns
Stun batons are an unusual modern variation designed to administer an electric shock in order to incapacitate the target. They consist of an insulated handle and guard, and a rigid shaft usually a foot or more in length for delivering a shock. Many designs function like an elongated stun gun or a cattle prod, requiring the tip to be held against the target and then manually triggering a shock by a switch in the handle. Some more sophisticated designs carry a charge along the shaft's entire surface, administering a shock on contact. This later design is especially useful in preventing the officer from having their weapon grabbed and taken away by an assailant.
Most batons of this design were not intended to be used as impact weapons and will break if used in this way, though a few were built to withstand occasional lighter impacts. They are rarely issued to patrol officers in modern times due to their price and the other associated problems with electroshock weapons.
Jitte
The jitte was a Japanese Edo period police weapon consisting of a round or octagonal metal rod about long with a hook-like guard above the handle. It was used in a similar manner to modern police batons and it continued to be issued in Japan to some police departments until the early 20th century. The jitte eventually inspired an early form of expandable baton called a tokushu keibo in the 1960s.
Improvised
Some non-purpose-built items have been used by law enforcement over the centuries as impact weapons. Some examples include:
Crowbars
Baseball bats
Oars or Paddles
Pickaxe handles
Flashlights
Although the Kel-Lite in the 1970s appears to have been the third flashlight designed specifically to be useful as an emergency weapon, the best-known example is the large, metal D-cell Maglite, still in use by some law enforcement and security personnel. Use of such flashlights as a club or baton is generally officially discouraged by the manufacturers and law enforcement officials, but its use is an option. As with all police weapons, there have been many examples of misuse, such as in the Malice Green beating in Detroit. The use of flashlights as improvised impact weapons is subject to the same use of force regulations as the use of purpose-designed impact weapons like batons.
Police officers may often choose to use such flashlights because they are viewed primarily as illumination devices; thus, if a police officer carries one in their hands during nighttime encounters with potentially violent subjects, it would be less likely to escalate the situation (by making the subject feel threatened) than if the officer were to be equipped with a baton or pepper spray canister instead. This permits the officer to appear less threatening while having an impact weapon in hand and ready for instantaneous action, should the situation indeed turn violent.
Characteristic of a flashlight used as a baton or club is the grip employed. Flashlights are commonly held with the bulb end pointing from the thumb side of the hand, such that it is pointing outward from the body when held palm upward. When wielded as a club, the bulb end points inward when the hand is palm upward, and the grip is closely choked to the bulb end. Another advantage to using a flashlight as a club is that in poorly lit situations it can be used to initially dazzle the eyes of an opponent. Law enforcement officers often deliberately shine flashlight beams into the eyes of suspects at night to cause temporary night-blindness as a preemptive defensive measure, whether or not the individual is likely to behave violently. The weight of a flashlight makes it a clumsy baton, unable to be swung swiftly.
Legality
Batons are legal for sworn law enforcement and military in most countries around the world. However, the legality of civilian carry for purpose-built batons varies greatly by country, and by local jurisdictions.
Brazil
There are no restrictions about batons to the general public, but private security guards can only carry wooden or rubber batons (no length is specified) according to Law 7102/83. They may also carry electric shock batons if they have a Less-Lethal Certification course. There is a general belief in Brazil that rubber batons are less prone to break bones than the wooden ones.
Canada
There is no law that prohibits batons; except for spring-loaded batons, which are defined as a prohibited weapon under a regulation entitled "Regulations Prescribing Certain Firearms and other Weapons, Components and Parts of Weapons, Accessories, Cartridge Magazines, Ammunition and Projectiles as Prohibited or Restricted" (also capable of being referred to by its registration number: SOR 98–462). However, it is a crime under section 90 of the Criminal Code to carry any weapon, including a baton, in a concealed fashion.
Hong Kong
According to Cap 217 (Weapons Ordinance), Laws of Hong Kong, any person who has possession of any prohibited items commits an offence, which includes expandable batons.
Indian subcontinent
In India and Bangladesh, police often carry a large bamboo stick called a lathi (, , ) that is used during riot controlling or used when a person is arrested or for self or public defence.
Ireland
In Ireland, telescopic truncheons are classified as illegal offensive weapons.
Sweden
All types of batons can be owned but not carried in public spaces by private citizens according to law (1988:254).
United Kingdom
Straight, side-handled (PR-24) and friction-lock batons were added to the list of offensive weapons in 2004 (except Scotland, where they were added in 2005), which prohibited their manufacture, sale, hire, offering for sale or hire, lending or giving to any other person under Section 141 Criminal Justice Act 1988. A loophole exists by way of martial arts weapons such as the Tonfa being legal to own, which is the exact same design as the PR-24 baton. The telescopic truncheon – defined as being a truncheon which extends automatically by hand pressure applied to a button, spring or other device in or attached to its handle – was banned in the original 1988 order.
Section 46 of the Offensive Weapons Act 2019, passed in May 2019, prohibits possession even in a private dwelling (e.g. home, closed off building site, behind a sales counter, etc.) previously, possession in private was permitted after meeting certain conditions based on ownership.
United States
Legality is determined by the laws of the individual states. Some, such as Vermont or Arizona, allow for legal carry in the absence of unlawful behavior or criminal intent. Others previously prohibited possession but constitutional challenges have overturned the bans, e.g. Connecticut v. DeCiccio (2009) and Hawaii. California has a general prohibition against the carrying of all "club" weapons by non-law enforcement. Constitutional challenges to California's law had failed prior to the United States Supreme Court's decision in New York State Rifle & Pistol Association, Inc. v. Bruen. However, in 2024, US District Court Judge Roger Benitez found that California's ban on club-like weapons was unconstitutional in light of the Supreme Court's Bruen opinion.
Jurisdictions with general prohibitions will sometimes make exceptions for persons employed as security guards or bodyguards, will provide for permits to be obtained for legal carry, or make exceptions for persons who complete an appropriate training course.
| Technology | Melee weapons | null |
12844718 | https://en.wikipedia.org/wiki/Dimorphic%20fungus | Dimorphic fungus | A dimorphic fungus is a fungus that can exist in the form of both mold and yeast. As this is usually brought about by a change in temperature, this fungus type is also described as a thermally dimorphic fungus. An example is Talaromyces marneffei, a human pathogen that grows as a mold at room temperature, and as a yeast at human body temperature.
The term dimorphic is commonly used for fungi that can grow both as yeast and filamentous cells, however many of these dimorphic fungi actually can grow in more than these two forms. Dimorphic is thus often used as a general reference for fungi being able to switch between yeast and filamentous cells, but not necessary limiting more shapes.
Ecology of dimorphic fungi
Several species of dimorphic fungi are important pathogens of humans and other animals, including Coccidioides immitis, Paracoccidioides brasiliensis, Candida albicans, Blastomyces dermatitidis, Histoplasma capsulatum, Sporothrix schenckii, and Emmonsia sp.
Some diseases caused by the fungi are:
sporotrichosis
blastomycosis
histoplasmosis
coccidioidomycosis
paracoccidioidomycosis
talaromycosis
candidiasis
Many other fungi, including the plant pathogen Ustilago maydis and the cheesemaker's fungus Geotrichum candidum also have dimorphic life cycles.
Mnemonics
In medical mycology, these memory aids help students remember that among human pathogens, dimorphism largely reflects temperature:
Mold in the Cold, Yeast in the Heat (Beast)
Body Heat Probably (Changes) Shape
Blastomyces dermatitidis, Histoplasma capsulatum, Paracoccidioides brasiliensis, (Coccidioides immitis) is in parentheses because it changes to a spherule of endospores, not yeast, in the heat), Sporothrix schenckii.
This phrase says "Probably" because there is always an exception (in this case fungi like Candida albicans) which change in the opposite direction: to mold in the heat!
| Biology and health sciences | Basics | Plants |
10592503 | https://en.wikipedia.org/wiki/Charging%20station | Charging station | A charging station, also known as a charge point, chargepoint, or electric vehicle supply equipment (EVSE), is a power supply device that supplies electrical power for recharging plug-in electric vehicles (including battery electric vehicles, electric trucks, electric buses, neighborhood electric vehicles, and plug-in hybrid vehicles).
There are two main types of EV chargers: Alternating current (AC) charging stations and direct current (DC) charging stations. Electric vehicle batteries can only be charged by direct current electricity, while most mains electricity is delivered from the power grid as alternating current. For this reason, most electric vehicles have a built-in AC-to-DC converter commonly known as the "onboard charger" (OBC). At an AC charging station, AC power from the grid is supplied to this onboard charger, which converts it into DC power to recharge the battery. DC chargers provide higher power charging (which requires much larger AC-to-DC converters) by building the converter into the charging station instead of the vehicle to avoid size and weight restrictions. The station then directly supplies DC power to the vehicle, bypassing the onboard converter. Most modern electric car models can accept both AC and DC power.
Charging stations provide connectors that conform to a variety of international standards. DC charging stations are commonly equipped with multiple connectors to charge various vehicles that use competing standards.
Public charging stations
Public charging stations are typically found street-side or at retail shopping centers, government facilities, and other parking areas. Private charging stations are usually found at residences, workplaces, and hotels.
Standards
Multiple standards have been established for charging technology to enable interoperability across vendors. Standards are available for nomenclature, power, and connectors. Tesla developed proprietary technology in these areas and began building its charging networking in 2012.
Nomenclature
In 2011, the European Automobile Manufacturers Association (ACEA) defined the following terms:
Socket outlet: the port on the electric vehicle supply equipment (EVSE) that supplies charging power to the vehicle
Plug: the end of the flexible cable that interfaces with the socket outlet on the EVSE. The socket outlet and plug are not used in North America because the cable is permanently attached.
Cable: a flexible bundle of conductors that connects the EVSE with the electric vehicle
Connector: the end of the flexible cable that interfaces with the vehicle inlet
Vehicle inlet: the port on the electric vehicle that receives charging power
The terms "electric vehicle connector" and "electric vehicle inlet" were previously defined in the same way under Article 625 of the United States National Electric Code (NEC) of 1999. NEC-1999 also defined the term "electric vehicle supply equipment" as the entire unit "installed specifically for the purpose of delivering energy from the premises wiring to the electric vehicle", including "conductors ... electric vehicle connectors, attachment plugs, and all other fittings, devices, power outlets, or apparatuses".
Tesla, Inc. uses the term charging station as the location of a group of chargers, and the term connector for an individual EVSE.
Voltage and power
Early standards
The National Electric Transportation Infrastructure Working Council (IWC) was formed in 1991 by the Electric Power Research Institute with members drawn from automotive manufacturers and the electric utilities to define standards in the United States; early work by the IWC led to the definition of three levels of charging in the 1999 National Electric Code (NEC) Handbook.
Under the 1999 NEC, Level 1 charging equipment (as defined in the NEC handbook but not in the code) was connected to the grid through a standard NEMA 5-20R 3-prong electrical outlet with grounding, and a ground-fault circuit interrupter was required within of the plug. The supply circuit required protection at 125% of the maximum rated current; for example, charging equipment rated at 16 amperes ("amps" or "A") continuous current required a breaker sized to 20 A.
Level 2 charging equipment (as defined in the handbook) was permanently wired and fastened at a fixed location under NEC-1999. It also required grounding and ground-fault protection; in addition, it required an interlock to prevent vehicle startup during charging and a safety breakaway for the cable and connector. A 40 A breaker (125% of continuous maximum supply current) was required to protect the branch circuit. For convenience and speedier charging, many early EVs preferred that owners and operators install Level 2 charging equipment, which was connected to the EV either through an inductive paddle (Magne Charge) or a conductive connector (Avcon).
Level 3 charging equipment used an off-vehicle rectifier to convert the input AC power to DC, which was then supplied to the vehicle. At the time it was written, the 1999 NEC handbook anticipated that Level 3 charging equipment would require utilities to upgrade their distribution systems and transformers.
SAE
The Society of Automotive Engineers (SAE International) defines the general physical, electrical, communication, and performance requirements for EV charging systems used in North America, as part of standard SAE J1772, initially developed in 2001. SAE J1772 defines four levels of charging, two levels each for AC and DC supplies; the differences between levels are based upon the power distribution type, standards and maximum power.
Alternating current (AC)
AC charging stations connect the vehicle's onboard charging circuitry directly to the AC supply.
AC Level 1: Connects directly to a standard 120V North American outlet; capable of supplying 616A (0.71.92kilowatts or "kW") depending on the capacity of a dedicated circuit.
AC Level 2: Uses 240V (single phase) or 208V (three phase) power to supply between 6 and 80A (1.419.2kW). It provides a significant charging speed increase over AC Level 1 charging.
Direct current (DC)
Commonly, though incorrectly, called "Level 3" charging based on the older NEC-1999 definition, DC charging is categorized separately in the SAE standard. In DC fast-charging, grid AC power is passed through an AC-to-DC converter in the station before reaching the vehicle's battery, bypassing any AC-to-DC converter on board the vehicle.
DC Level 1: Supplies a maximum of 80kW at 501000V.
DC Level 2: Supplies a maximum of 400kW at 501000V.
Additional standards released by SAE for charging include SAE J3068 (three-phase AC charging, using the Type 2 connector defined in IEC 62196-2) and SAE J3105 (automated connection of DC charging devices).
IEC
In 2003, the International Electrotechnical Commission (IEC) adopted a majority of the SAE J1772 standard under IEC 62196-1 for international implementation.
The IEC alternatively defines charging in modes (IEC 61851-1):
Mode 1: slow charging from a regular electrical socket (single- or three-phase AC)
Mode 2: slow charging from a regular AC socket but with some EV-specific protection arrangement (i.e. the Park & Charge or the PARVE systems)
Mode 3: slow or fast AC charging using a specific EV multi-pin socket with control and protection functions (i.e. SAE J1772 and IEC 62196-2)
Mode 4: DC fast charging using a specific charging interface (i.e. IEC 62196-3, such as CHAdeMO)
The connection between the electric grid and "charger" (electric vehicle supply equipment) is defined by three cases (IEC 61851-1):
Case A: any charger connected to the mains (the mains supply cable is usually attached to the charger) usually associated with modes 1 or 2.
Case B: an on-board vehicle charger with a mains supply cable that can be detached from both the supply and the vehicle – usually mode 3.
Case C: DC dedicated charging station. The mains supply cable may be permanently attached to the charge station as in mode 4.
Tesla NACS
The North American Charging System (NACS) was developed by Tesla, Inc. for use in the company's vehicles. It remained a proprietary standard until 2022 when its specifications were published by Tesla. The connector is physically smaller than the J1172/CCS connector, and uses the same pins for both AC and DC charging functionality.
As of November 2023, automakers Ford, General Motors, Rivian, Volvo, Polestar, Mercedes-Benz, Nissan, Honda, Jaguar, Fisker, Hyundai, BMW, Toyota, Subaru, and Lucid Motors have all committed to equipping their North American vehicles with NACS connectors in the future. Automotive startup Aptera Motors has also adopted the connector standard in its vehicles. Other automakers, such as Stellantis and Volkswagen have not made an announcement.
To meet European Union (EU) requirements on recharging points, Tesla vehicles sold in the EU are equipped with a CCS Combo 2 port. Both the North America and the EU port take 480V DC fast charging through Tesla's network of Superchargers, which variously use NACS and CCS charging connectors. Depending on the Supercharger version, power is supplied at 72, 150, or 250 kW, the first corresponding to DC Level 1 and the second and third corresponding to DC Level 2 of SAE J1772. As of Q4 2021, Tesla reported 3,476 supercharging locations worldwide and 31,498 supercharging chargers (about 9 chargers per location on average).
Future development
An extension to the CCS DC fast-charging standard for electric cars and light trucks is under development, which will provide higher power charging for large commercial vehicles (Class 8, and possibly 6 and 7 as well, including school and transit buses). When the Charging Interface Initiative e. V. (CharIN) task force was formed in March 2018, the new standard being developed was originally called High Power Charging (HPC) for Commercial Vehicles (HPCCV), later renamed Megawatt Charging System (MCS). MCS is expected to operate in the range of 2001500V and 03000A for a theoretical maximum power of 4.5megawatts (MW). The proposal calls for MCS charge ports to be compatible with existing CCS and HPC chargers. The task force released aggregated requirements in February 2019, which called for maximum limits of 1000V DC (optionally, 1500V DC) and 3000A continuous rating.
A connector design was selected in May 2019 and tested at the National Renewable Energy Laboratory (NREL) in September 2020. Thirteen manufacturers participated in the test, which checked the coupling and thermal performance of seven vehicle inlets and eleven charger connectors. The final connector requirements and specification was adopted in December 2021 as MCS connector version 3.2.
With support from Portland General Electric, on 21 April 2021 Daimler Trucks North America opened the "Electric Island", the first heavy-duty vehicle charging station, across the street from its headquarters in Portland, Oregon. The station is capable of charging eight vehicles simultaneously, and the charging bays are sized to accommodate tractor-trailers. In addition, the design is capable of accommodating >1MW chargers once they are available. A startup company, WattEV, announced plans in May 2021 to build a 40-stall truck stop/charging station in Bakersfield, California. At full capacity, it would provide a combined 25MW of charging power, partially drawn from an on-site solar array and battery storage.
Connectors
Common connectors include Type 1 (Yazaki), Type 2 (Mennekes), CCS Combo 1 and 2, CHAdeMO, and Tesla. Many standard plug types are defined in IEC 62196-2 (for AC supplied power) and 62196-3 (for DC supplied power):
Type 1: single-phase AC vehicle coupler – SAE J1772/2009 automotive plug specifications
Type 2: single- and three-phase AC vehicle coupler – VDE-AR-E 2623-2-2, SAE J3068, and GB/T 20234.2 plug specifications
Type 3: single- and three-phase AC vehicle coupler equipped with safety shutters – EV Plug Alliance proposal
Type 4: DC fast charge couplers
Configuration AA: CHAdeMO
Configuration BB: GB/T 20234.3
Configurations CC/DD: (reserved)
Configuration EE: CCS Combo 1
Configuration FF: CCS Combo 2
| Technology | Concepts of ground transport | null |
19637796 | https://en.wikipedia.org/wiki/Scarecrow | Scarecrow | A scarecrow is a decoy or mannequin that is often in the shape of a human. Humanoid scarecrows are usually dressed in old clothes and placed in open fields to discourage birds from disturbing and feeding on recently cast seed and growing crops. Scarecrows are used around the world by farmers, and are a notable symbol of farms and the countryside in popular culture.
Design
The common form of a scarecrow is a humanoid figure dressed in old clothes and placed in open fields to discourage birds such as crows or sparrows from disturbing and feeding on recently cast seed and growing crops. Machinery such as windmills have been employed as scarecrows, but the effectiveness lessens as animals become familiar with the structures.
Since the invention of the humanoid scarecrow, more effective methods have been developed. On California farmland, highly-reflective aluminized PET film ribbons are tied to the plants to produce shimmers from the sun. Another approach is using automatic noise guns powered by propane gas. One winery in New York has even used inflatable tube men or airdancers to scare away birds.
Cultural impact
Joe's Scarecrow Village in Cape Breton, Canada, was a roadside attraction displaying dozens of scarecrows.
The Japanese village of Nagoro, on the island of Shikoku in the Tokushima Prefecture, has 35 inhabitants but more than 350 scarecrows.
In the United Kingdom, where there are a few different languages and several different dialects, there are a wide range of alternative names such as:
A supervillain named Scarecrow is part of the Rogues Gallery of the DC hero Batman.
Festivals
In England, the Urchfont Scarecrow Festival was established in the 1990s and has become a major local event, attracting up to 10,000 people annually for the May Day Bank Holiday. Originally based on an idea imported from Derbyshire, or Kettlewell, North Yorkshire, it was the first Scarecrow Festival to be established in the whole of southern England.
Belbroughton, north Worcestershire, holds an annual Scarecrow Weekend on the last weekend of each September since 1996, which raises money for local charities. The village of Meerbrook in Staffordshire holds an annual Scarecrow Festival during the month of May. Tetford and Salmonby, Lincolnshire, jointly host one.
The festival at Wray, Lancashire, was established in the early 1990s and continues to the present day. In the village of Orton, Eden, Cumbria scarecrows are displayed each year, often using topical themes such as a Dalek exterminating a Wind turbine to represent local opposition to a wind farm.
The village of Blackrod, near Bolton in Greater Manchester, holds a popular annual Scarecrow Festival over a weekend usually in early July.
Norland, West Yorkshire, has a Scarecrow festival. Kettlewell in North Yorkshire has held an annual festival since 1994. The villages of Cotherstone, Staindrop, and Middleton-in-Teesdale in County Durham have annual scarecrow festivals.
Scotland's first scarecrow festival was held in West Kilbride, North Ayrshire, in 2004, and there is also one held in Montrose. On the Isle of Skye, the Tattie bogal event is held each year, featuring a scarecrow trail and other events. Tonbridge in Kent also host an annual Scarecrow Trail, organised by the local Rotary Club to raise money for local charities. Gisburn, Lancashire, held its first Scarecrow Festival in June 2014.
Mullion, in Cornwall, has an annual scarecrow festival since 2007.
In the US, St. Charles, Illinois, hosts an annual Scarecrow Festival. Peddler's Village in Bucks County, Pennsylvania, hosts an annual scarecrow festival and presents a scarecrow display in September–October that draws tens of thousands of visitors.
The "pumpkin people" come in the autumn months in the valley region of Nova Scotia, Canada. They are scarecrows with pumpkin heads applied to them doing various things such as playing the fiddle or riding a wooden horse. Hickling, in the south of Nottinghamshire, is another village that celebrates an annual scarecrow event. It is very popular and has successfully raised a great deal of money for charity. Meaford, Ontario, has celebrated the Scarecrow Invasion since 1996.
In the Philippines in 2015, the Province of Isabela started a scarecrow festival named after the local language: the Bambanti Festival. The province invites all its cities and towns to participate for the festivities, which last a week; it has drawn tourists from around the island of Luzon.
The largest gathering of scarecrows in one location is 3,812 and was achieved by National Forest Adventure Farm in Burton-upon-Trent, Staffordshire, UK, on 7 August 2014.
Gallery
| Technology | Pest and disease control | null |
19637950 | https://en.wikipedia.org/wiki/Housefly | Housefly | The housefly (Musca domestica) is a fly of the suborder Cyclorrhapha. It possibly originated in the Middle East, and spread around the world as a commensal of humans. It is the most common fly species found in houses. Adults are gray to black, with four dark, longitudinal lines on the thorax, slightly hairy bodies, and a single pair of membranous wings. They have red eyes, set farther apart in the slightly larger female.
The female housefly usually mates only once and stores the sperm for later use. It lays batches of about 100 eggs on decaying organic matter such as food waste, carrion, or feces. These soon hatch into legless white larvae, known as maggots. After two to five days of development, these metamorphose into reddish-brown pupae, about long. Adult flies normally live for two to four weeks, but can hibernate during the winter. The adults feed on a variety of liquid or semi-liquid substances, as well as solid materials which have been softened by their saliva. They can carry pathogens on their bodies and in their feces, contaminate food, and contribute to the transfer of food-borne illnesses, while, in numbers, they can be physically annoying. For these reasons, they are considered pests.
Houseflies, with short life cycles and ease with which they can be maintained, have been found useful for laboratory research into aging and sex determination. Houseflies appear in literature from Ancient Greek myth and Aesop's "The Impertinent Insect" onwards. Authors sometimes choose the housefly to speak of the brevity of life, as in William Blake's 1794 poem "The Fly", which deals with mortality subject to uncontrollable circumstances.
Description
Adult houseflies are usually long with a wingspan of . The females tend to be larger winged than males, while males have relatively longer legs. Females tend to vary more in size and there is geographic variation with larger individuals in higher latitudes. The head is strongly convex in front and flat and slightly conical behind. The pair of large compound eyes almost touch in the male, but are more widely separated in the female. They have three simple eyes (ocelli) and a pair of short antennae. Houseflies process visual information around seven times more quickly than humans, enabling them to identify and avoid attempts to catch or swat them, since they effectively see the human's movements in slow motion with their higher flicker fusion rate.
The mouthparts are specially adapted for a liquid diet; the mandibles and maxillae are reduced and not functional, and the other mouthparts form a retractable, flexible proboscis with an enlarged, fleshy tip, the labellum. This is a sponge-like structure that is characterized by many grooves, called pseudotracheae, which suck up fluids by capillary action. It is also used to distribute saliva to soften solid foods or collect loose particles. Houseflies have chemoreceptors, organs of taste, on the tarsi of their legs, so they can identify foods such as sugars by walking over them. Houseflies are often seen cleaning their legs by rubbing them together, enabling the chemoreceptors to taste afresh whatever they walk on next. At the end of each leg is a pair of claws, and below them are two adhesive pads, pulvilli, enabling the housefly to walk up smooth walls and ceilings using Van der Waals forces. The claws help the housefly to unstick the foot for the next step. Houseflies walk with a common gait on horizontal and vertical surfaces with three legs in contact with the surface and three in movement. On inverted surfaces, they alter the gait to keep four feet stuck to the surface. Houseflies land on a ceiling by flying straight towards it; just before landing, they make a half roll and point all six legs at the surface, absorbing the shock with the front legs and sticking a moment later with the other four.
The thorax is a shade of gray, sometimes even black, with four dark, longitudinal bands of even width on the dorsal surface. The whole body is covered with short hairs. Like other Diptera, houseflies have only one pair of wings; what would be the hind pair is reduced to small halteres that aid in flight stability. The wings are translucent with a yellowish tinge at their base. Characteristically, the medial vein (M1+2 or fourth long vein) shows a sharp upward bend. Each wing has a lobe at the back, the calypter, covering the haltere. The abdomen is gray or yellowish with a dark stripe and irregular dark markings at the side. It has 10 segments which bear spiracles for respiration. In males, the ninth segment bears a pair of claspers for copulation, and the 10th bears anal cerci in both sexes.
A variety of species around the world appear similar to the housefly, such as the lesser house fly, Fannia canicularis; the stable fly, Stomoxys calcitrans; and other members of the genus Musca such as M. vetustissima, the Australian bush fly and several closely related taxa that include M. primitiva, M. shanghaiensis, M. violacea, and M. varensis. The systematic identification of species may require the use of region-specific taxonomic keys and can require dissections of the male reproductive parts for confirmation.
Distribution
The housefly is probably the insect with the widest distribution in the world; it is largely associated with humans and has accompanied them around the globe. It is present in the Arctic, as well as in the tropics, where it is abundant. It is present in all populated parts of Europe, Asia, Africa, Australasia, and the Americas.
Evolution and taxonomy
Though the order of flies (Diptera) is much older, true houseflies are believed to have evolved in the beginning of the Cenozoic Era. The housefly's superfamily, Muscoidea, is most closely related to the Oestroidea (blow flies, flesh flies and allies), and more distantly to the Hippoboscoidea (louse flies, bat flies and allies). They are thought to have originated in the southern Palearctic region, particularly the Middle East. Because of their close, commensal relationship with humans, they probably owe their worldwide dispersal to co-migration with humans.
The housefly was first described as Musca domestica in 1758 based on the common European specimens by the Swedish botanist and zoologist Carl Linnaeus in his Systema naturae and continues to be classified under that name. A more detailed description was given in 1776 by the Danish entomologist Johan Christian Fabricius in his Genera Insectorum.
Life cycle
Each female housefly can lay up to 500 eggs in her lifetime, in several batches of about 75 to 150. The eggs are white and are about in length, and they are deposited by the fly in a suitable place, usually dead and decaying organic matter, such as food waste, carrion, or feces. Within a day, larvae (maggots) hatch from the eggs; they live and feed where they were laid. They are pale-whitish, long, thinner at the mouth end, and legless. Larval development takes from two weeks, under optimal conditions, to 30 days or more in cooler conditions. The larvae avoid light; the interiors of heaps of animal manure provide nutrient-rich sites and ideal growing conditions, warm, moist, and dark.
At the end of their third instar, the larvae crawl to a dry, cool place and transform into pupae. The pupal case is cylindrical with rounded ends, about long, and formed from the last shed larval skin. It is yellowish at first, darkening through red and brown to nearly black as it ages. Pupae complete their development in two to six days at , but may take 20 days or more at .
When metamorphosis is complete, the adult housefly emerges from the pupa. To do this, it uses the ptilinum, an eversible pouch on its head, to tear open the end of the pupal case. Having emerged from the pupa, it ceases to grow; a small fly is not necessarily a young fly, but is instead the result of getting insufficient food during the larval stage.
Male houseflies are sexually mature after 16 hours and females after 24. Females produce a pheromone, (Z)-9-tricosene (muscalure). This cuticular hydrocarbon is not released into the air and males sense it only on contact with females; it has found use as in pest control, for luring males to fly traps. The male initiates the mating by bumping into the female, in the air or on the ground, known as a "strike". He climbs on to her thorax, and if she is receptive, a courtship period follows, in which the female vibrates her wings and the male strokes her head. The male then reverses onto her abdomen and the female pushes her ovipositor into his genital opening; copulation, with sperm transfer, lasts for several minutes. Females normally mate only once and then reject further advances from males, while males mate multiple times. A volatile semiochemical that is deposited by females on their eggs attracts other gravid females and leads to clustered egg deposition.
The larvae depend on warmth and sufficient moisture to develop; generally, the warmer the temperature, the faster they grow. In general, fresh swine and chicken manures present the best conditions for the developing larvae, reducing the larval period and increasing the size of the pupae. Cattle, goat, and horse manures produce fewer, smaller pupae, while mature swine manure composted with water content under 30%, approached 100% mortality of the larvae. Pupae can range from about in weight under different conditions.
The life cycle can be completed in seven to ten days under optimal conditions, but may take up to two months in adverse circumstances. In temperate regions, 12 generations may occur per year, and in the tropics and subtropics, more than 20.
Ecology
Houseflies play an important ecological role in breaking down and recycling organic matter. Adults are mainly carnivorous; their primary food is animal matter, carrion, and feces, but they also consume milk, sugary substances, and rotting fruit and vegetables. Solid foods are softened with saliva before being sucked up. They can be opportunistic blood feeders. Houseflies have a mutualistic relationship with the bacterium Klebsiella oxytoca, which can live on the surface of housefly eggs and deter fungi which compete with the housefly larvae for nutrients.
Adult houseflies are diurnal and rest at night. If inside a building after dark, they tend to congregate on ceilings, beams, and overhead wires, while out of doors, they crawl into foliage or long grass, or rest in shrubs and trees or on wires. In cooler climates, some houseflies hibernate in winter, choosing to do so in cracks and crevices, gaps in woodwork, and the folds of curtains. They arouse in the spring when the weather warms up, and search out a place to lay their eggs.
Houseflies have many predators, including birds, reptiles, amphibians, various insects, and spiders. The eggs, larvae, and pupae have many species of stage-specific parasites and parasitoids. Some of the more important are the parasitic wasps Muscidifurax uniraptor and Spalangia cameroni; these lay their eggs in the housefly larvae tissue and their offspring complete their development before the adult houseflies can emerge from the pupae. Hister beetles feed on housefly larvae in manure heaps and the predatory mite Macrocheles muscae domesticae consumes housefly eggs, each mite eating 20 eggs per day.
Houseflies sometimes carry phoretic (nonparasitic) passengers, including mites such as Macrocheles muscaedomesticae and the pseudoscorpion Lamprochernes chyzeri.
The pathogenic fungus Entomophthora muscae causes a fatal disease in houseflies. After infection, the fungal hyphae grow throughout the body, killing the housefly in about five days. Infected houseflies have been known to seek high temperatures that could suppress the growth of the fungus. Affected females tend to be more attractive to males, but the fungus-host interactions have not been fully understood. The housefly also acts as the alternative host to the parasitic nematode Habronema muscae that attacks horses. A virus that causes enlargement of the salivary glands, salivary gland hypertrophy virus (SGHV), is spread among houseflies through contact with food and infected female houseflies become sterile.
Relationship with humans
Houseflies are often considered a nuisance, disturbing people while at leisure and at work, but they are disliked principally because of their habits of contaminating foodstuffs. They alternate between breeding and feeding in dirty places with feeding on human foods, during which process they soften the food with saliva and deposit their feces, creating a health hazard. However, housefly larvae are as nutritious as fish meal, and could be used to convert waste to insect-based animal feed for farmed fish and livestock. Housefly larvae have been used in traditional cures since the Ming period in China (1386 AD) for a range of medical conditions and have been considered as a useful source of chitosan, with antioxidant properties, and possibly other proteins and polysaccharides of medical value.
Houseflies have been used in art and artifacts in many cultures. In 16th- and 17th-century European vanitas paintings, houseflies sometimes occur as memento mori. They may also be used for other effects as in the Flemish painting, the Master of Frankfurt (1496). Housefly amulets were popular in ancient Egypt.
As a disease vector
Houseflies can fly for several kilometers from their breeding places, carrying a wide variety of organisms on their hairs, mouthparts, vomitus, and feces. Parasites carried include cysts of protozoa, e.g. Entamoeba histolytica and Giardia lamblia and eggs of helminths; e.g., Ascaris lumbricoides, Trichuris trichiura, Hymenolepis nana, and Enterobius vermicularis. Houseflies do not serve as a secondary host or act as a reservoir of any bacteria of medical or veterinary importance, but they do serve as mechanical vectors to over 100 pathogens, such as those causing typhoid, cholera, salmonellosis, bacillary dysentery, tuberculosis, anthrax, ophthalmia, and pyogenic cocci, making them especially problematic in hospitals and during outbreaks of certain diseases. Disease-causing organisms on the outer surface of the housefly may survive for a few hours, but those in the crop or gut can be viable for several days. Usually, too few bacteria are on the external surface of the houseflies (except perhaps for Shigella) to cause infection, so the main routes to human infection are through the housefly's regurgitation and defecation. A number of bacterial endosymbionts have however been detected in sequence-based identification from whole genome sequences extracted from flies, the greatest numbers being detected in the abdomen.
In the early 20th century, Canadian public health workers believed that the control of houseflies was important in controlling the spread of tuberculosis. A "swat that fly" contest was held for children in Montreal in 1912. Houseflies were targeted in 1916, when a polio epidemic broke out in the eastern United States. The belief that housefly control was the key to disease control continued, with extensive use of insecticidal spraying well until the mid-1950s, declining only after the introduction of Salk's vaccine. In China, Mao Zedong's Four Pests Campaign between 1958 and 1962 exhorted the people to catch and kill houseflies, along with rats, mosquitoes, and sparrows.
In warfare
During the Second World War, the Japanese worked on entomological warfare techniques under Shirō Ishii. Japanese Yagi bombs developed at Pingfan consisted of two compartments, one with houseflies and another with a bacterial slurry that coated the houseflies prior to release. Vibrio cholerae, which causes cholera, was the bacterium of choice, and was used by Japan against the Chinese in Baoshan in 1942, and in northern Shandong in 1943. The Baoshan bombing produced epidemics that killed 60,000 people in the initial stages, reaching a radius of which finally took a toll of 200,000 victims. The Shandong attack killed 210,000; the occupying Japanese troops had been vaccinated in advance.
In waste management
The ability of housefly larvae to feed and develop in a wide range of decaying organic matter is important for recycling of nutrients in nature. This could be exploited to combat ever-increasing amounts of waste. Housefly larvae can be mass-reared in a controlled manner in animal manure, reducing the bulk of waste and minimizing environmental risks of its disposal. Harvested maggots may be used as feed for animal nutrition.
Control
Houseflies can be controlled, at least to some extent, by physical, chemical, or biological means. Physical controls include screening with small mesh or the use of vertical strips of plastic or strings of beads in doorways to prevent entry of houseflies into buildings. Fans to create air movement or air barriers in doorways can deter houseflies from entering, and food premises often use fly-killing devices; sticky fly papers hanging from the ceiling are effective, but electric "bug zappers" should not be used directly above food-handling areas because of scattering of contaminated insect parts. Another approach is the elimination as far as possible of potential breeding sites. Keeping garbage in lidded containers and collecting it regularly and frequently, prevents any eggs laid from developing into adults. Unhygienic rubbish tips are a prime housefly-breeding site, but if garbage is covered by a layer of soil, preferably daily, this can be avoided.
Insecticides can be used. Larvicides kill the developing larvae, but large quantities may need to be used to reach areas below the surface. Aerosols can be used in buildings to "zap" houseflies, but outside applications are only temporarily effective. Residual sprays on walls or resting sites have a longer-lasting effect. Many strains of housefly have become immune to the most commonly used insecticides. Resistance to carbamates and organophosphates is conferred by variation in acetylcholinesterase genes. M. domestica has achieved a high degree of resistance. Resistance monitoring is vital to avoid continued use of ineffective active ingredients such as found in the notably severe example of Freeman et al 2019 in Kansas and Maryland, USA.
Several means of biological pest control have been investigated. These include the introduction of another species, the black soldier fly (Hermetia illucens), whose larvae compete with those of the housefly for resources. The introduction of dung beetles to churn up the surface of a manure heap and render it unsuitable for breeding is another approach. Augmentative biological control by releasing parasitoids can be used, but houseflies breed so fast that the natural enemies are unable to keep up.
In science
The ease of culturing houseflies, and the relative ease of handling them when compared to the fruit fly Drosophila, have made them useful as model organism for use in laboratories. The American entomologist Vincent Dethier, in his humorous To Know A Fly (1962), pointed out that as a laboratory animal, houseflies did not trouble anyone sensitive to animal cruelty. Houseflies have a small number of chromosomes, haploid 6 or diploid 12. Because the somatic tissue of the housefly consists of long-lived postmitotic cells, it can be used as an informative model system for understanding cumulative age-related cellular alterations. Oxidative DNA damage 8-hydroxydeoxyguanosine in houseflies was found in one study to increase with age and reduce life expectancy supporting the hypothesis that oxidative molecular damage is a causal factor in senescence (aging).
The housefly is an object of biological research, partly for its variable sex-determination mechanism. Although a wide variety of sex-determination mechanisms exists in nature (e.g. male and female heterogamy, haplodiploidy, environmental factors), the way sex is determined is usually fixed within a species. The housefly is, however, thought to exhibit multiple mechanisms for sex determination, such as male heterogamy (like most insects and mammals), female heterogamy (like birds), and maternal control over offspring sex. This is because a male-determining gene (Mdmd) can be found on most or all housefly chromosomes. Sexual differentiation is controlled, as in other insects, by an ancient developmental switch, doublesex, which is regulated by the transformer protein in many different insects. Mdmd causes male development by negatively regulating transformer. There is also a female-determining allele of transformer that is not sensitive to the negative regulation of Mdmd.
The antimicrobial peptides produced by housefly maggots are of pharmacological interest.
In the 1970s, the aircraft modeler Frank Ehling constructed miniature balsa-wood aircraft powered by live houseflies. Studies of tethered houseflies have helped in the understanding of insect vision, sensory perception, and flight control.
In literature
The Impertinent Insect is a group of five fables, sometimes ascribed to Aesop, concerning an insect, in one version a fly, which puffs itself up to seem important. In the Biblical fourth plague of Egypt, flies represent death and decay, while the Philistine god Beelzebub's name may mean "lord of the flies". In Greek mythology, Myiagros was a god who chased away flies during the sacrifices to Zeus and Athena; Zeus sent a fly to bite Pegasus, causing Bellerophon to fall back to Earth when he attempted to ride the winged steed to Mount Olympus. In the traditional Navajo religion, Big Fly is an important spirit being.
William Blake's 1794 poem "The Fly", part of his collection Songs of Experience, deals with the insect's mortality, subject to uncontrollable circumstances, just like humans. Emily Dickinson's 1855 poem "I Heard a Fly Buzz When I Died" speaks of flies in the context of death. In William Golding's 1954 novel Lord of the Flies, the fly is, however, a symbol of the children involved.
Ogden Nash's humorous two-line 1942 poem "God in His wisdom made the fly/And then forgot to tell us why." indicates the debate about the value of biodiversity, given that even those considered by humans as pests have their place in the world's ecosystems.
| Biology and health sciences | Flies (Diptera) | null |
19637986 | https://en.wikipedia.org/wiki/Gingivitis | Gingivitis | Gingivitis is a non-destructive disease that causes inflammation of the gums; ulitis is an alternative term. The most common form of gingivitis, and the most common form of periodontal disease overall, is in response to bacterial biofilms (also called plaque) that are attached to tooth surfaces, termed plaque-induced gingivitis. Most forms of gingivitis are plaque-induced.
While some cases of gingivitis never progress to periodontitis, periodontitis is always preceded by gingivitis.
Gingivitis is reversible with good oral hygiene; however, without treatment, gingivitis can progress to periodontitis, in which the inflammation of the gums results in tissue destruction and bone resorption around the teeth. Periodontitis can ultimately lead to tooth loss.
Signs and symptoms
The symptoms of gingivitis are somewhat non-specific and manifest in the gum tissue as the classic signs of inflammation:
Swollen gums
Bright red gums
Gums that are tender or painful to the touch
Bleeding gums or bleeding after brushing and/or flossing
Bad breath (halitosis)
Additionally, the stippling that normally exists in the gum tissue of some individuals will often disappear and the gums may appear shiny when the gum tissue becomes swollen and stretched over the inflamed underlying connective tissue. The accumulation may also emit an unpleasant odor. When the gingiva are swollen, the epithelial lining of the gingival crevice becomes ulcerated and the gums will bleed more easily with even gentle brushing, and especially when flossing.
Complications
Recurrence of gingivitis
Periodontitis
Infection or abscess of the gingiva or the jaw bones
Trench mouth (bacterial infection and ulceration of the gums)
Swollen lymph nodes
Associated with premature birth and low birth weight
Alzheimer's and dementia
A study from 2018 found evidence that gingivitis bacteria may be linked to Alzheimer's disease. Scientists agree that more research is needed to prove a cause and effect link. "Studies have also found that the bacteria P. gingivalis – which are responsible for many forms of gum disease – can migrate from the mouth to the brain in mice. And on entry to the brain, P. gingivalis can reproduce all of the characteristic features of Alzheimer's disease."
Cause
The cause of plaque-induced gingivitis is bacterial plaque, which acts to initiate the body's host response. This, in turn, can lead to destruction of the gingival tissues, which may progress to destruction of the periodontal attachment apparatus. The plaque accumulates in the small gaps between teeth, in the gingival grooves and in areas known as plaque traps: locations that serve to accumulate and maintain plaque. Examples of plaque traps include bulky and overhanging restorative margins, clasps of removable partial dentures and calculus (tartar) that forms on teeth. Although these accumulations may be tiny, the bacteria in them produce chemicals, such as degradative enzymes, and toxins, such as lipopolysaccharide (LPS, otherwise known as endotoxin) or lipoteichoic acid (LTA), that promote an inflammatory response in the gum tissue. This inflammation can cause an enlargement of the gingiva and subsequent formation. Early plaque in health consists of a relatively simple bacterial community dominated by Gram-positive cocci and rods. As plaque matures and gingivitis develops, the communities become increasingly complex with higher proportions of Gram-negative rods, fusiforms, filaments, spirilla and spirochetes. Later experimental gingivitis studies, using culture, provided more information regarding the specific bacterial species present in plaque. Taxa associated with gingivitis included Fusobacterium nucleatum subspecies polymorphum, Lachnospiraceae [G-2] species HOT100, Lautropia species HOTA94, and Prevotella oulorum (a species of Prevotella bacterium), whilst Rothia dentocariosa was associated with periodontal health. Further study of these taxa is warranted and may lead to new therapeutic approaches to prevent periodontal disease including systemic health.
Risk factors
Risk factors associated with gingivitis include the following:
age
osteoporosis
low dental care utilization
poor oral hygiene
overly aggressive oral hygiene such as brushing with stiff bristles
mouth breathing during sleep
orthodontic braces
medications and conditions that dry the mouth
cigarette smoking
genetic factors
stress
mental health issues such as depression
pre-existing conditions such as diabetes
drinking
obesity.
Diagnosis
Gingivitis is a category of periodontal disease in which there is no loss of bone but inflammation and bleeding are present.
Each tooth is divided into four gingival units (mesial, distal, buccal, and lingual) and given a score from 0–3 based on the gingival index. The four scores are then averaged to give each tooth a single score.
The diagnosis of the periodontal disease gingivitis is done by a dentist. The diagnosis is based on clinical assessment data acquired during a comprehensive periodontal exam. Either a registered dental hygienist or a dentist may perform the comprehensive periodontal exam but the data interpretation and diagnosis are done by the dentist. The comprehensive periodontal exam consists of a visual exam, a series of radiographs, probing of the gingiva, determining the extent of current or past damage to the periodontium and a comprehensive review of the medical and dental histories.
Current research shows that activity levels of the following enzymes in saliva samples are associated with periodontal destruction: aspartate aminotransferase (AST), alanine aminotransferase (ALT), gamma glutamyl transferase (GGT), alkaline phosphatase (ALP), and acid phosphatase (ACP). Therefore, these enzyme biomarkers may be used to aid in the diagnosis and treatment of gingivitis and periodontitis.
A dental hygienist or dentist will check for the symptoms of gingivitis, and may also examine the amount of plaque in the oral cavity. A dental hygienist or dentist will also look for signs of periodontitis using X-rays or periodontal probing as well as other methods.
If gingivitis is not responsive to treatment, referral to a periodontist (a specialist in diseases of the gingiva and bone around teeth and dental implants) for further treatment may be necessary.
Classification
1999 Classification
As defined by the 1999 World Workshop in Clinical Periodontics, there are two primary categories of gingival diseases, each with numerous subgroups:
Dental plaque-induced gingival diseases.
Gingivitis associated with plaque only
Gingival diseases modified by systemic factors
Gingival diseases modified by medications
Gingival diseases modified by malnutrition
Non-plaque-induced gingival lesions
Gingival diseases of specific bacterial origin
Gingival diseases of viral origin
Gingival diseases of fungal origin
Gingival diseases of genetic origin
Gingival manifestations of systemic conditions
Traumatic lesions
Foreign body reactions
Not otherwise specified
2017 Classification
As defined by the 2017 World Workshop, periodontal health, gingival diseases/ conditions have been categorised into the following:
Periodontal health and gingival health
Clinical gingival health on an intact periodontium
Clinical gingival health on a reduced periodontium
Stable periodontitis patient
Non-periodontitis patient
Gingivitis – dental biofilm-induced
Associated with dental biofilm alone
Mediated by systemic or local risk factors
Drug-influenced gingival enlargement
Gingival diseases – non-dental biofilm induced
Genetic/ developmental disorders
Specific infections
Inflammatory and immune conditions
Reactive processes
Neoplasms
Endocrine, nutritional & metabolic diseases
Traumatic lesions
Gingival pigmentation
Prevention
Gingivitis can be prevented through regular oral hygiene that includes daily brushing and flossing. Hydrogen peroxide, saline, alcohol or chlorhexidine mouth washes may also be employed. In a 2004 clinical study, the beneficial effect of hydrogen peroxide on gingivitis has been highlighted. The use of oscillation type brushes might reduce the risk of gingivitis compared to manual brushing.
Rigorous plaque control programs along with periodontal scaling and curettage also have proved to be helpful, although according to the American Dental Association, periodontal scaling and root planing are considered as a treatment for periodontal disease, not as a preventive treatment for periodontal disease. In a 1997 review of effectiveness data, the U.S. Food and Drug Administration (FDA) found clear evidence showing that toothpaste containing triclosan was effective in preventing gingivitis. In 2017 the FDA banned triclosan in many consumer products but allowed it to remain in toothpaste because of its effectiveness against gingivitis. In 2019, Colgate, under pressure from health advocates, removed triclosan from the last toothpaste on the market containing it, Colgate Total.
Treatment
The focus of treatment is to remove plaque. Therapy is aimed at the reduction of oral bacteria and may take the form of regular periodic visits to a dental professional together with adequate oral hygiene home care. Thus, several of the methods used in the prevention of gingivitis can also be used for the treatment of manifest gingivitis, such as scaling, root planing, curettage, mouth washes containing chlorhexidine or hydrogen peroxide, and flossing. Interdental brushes also help remove any causative agents.
Powered toothbrushes work better than manual toothbrushes in reducing the disease.
The active ingredients that "reduce plaque and demonstrate effective reduction of gingival inflammation over a period of time" are triclosan, chlorhexidine digluconate, and a combination of thymol, menthol, eucalyptol, and methyl salicylate. These ingredients are found in toothpaste and mouthwash. Hydrogen peroxide was long considered a suitable over-the-counter agent to treat gingivitis. There has been evidence to show the positive effect on controlling gingivitis in short-term use. A study indicates the fluoridated hydrogen peroxide-based mouth rinse can remove teeth stain and reduce gingivitis.
Based on a limited evidence, mouthwashes with essential oils may also be useful, as they contain ingredients with anti-inflammatory properties, such as thymol, menthol and eucalyptol.
The bacteria that causes gingivitis can be controlled by using an oral irrigator daily with a mouthwash containing an antibiotic. Either amoxicillin, cephalexin, or minocycline in 500 grams of a non-alcoholic fluoride mouthwash is an effective mixture.
Overall, intensive oral hygiene care has been shown to improve gingival health in individuals with well-controlled type 2 diabetes. Periodontal destruction is also slowed due to the extensive oral care. Intensive oral hygiene care (oral health education plus supra-gingival scaling) without any periodontal therapy improves gingival health, and may prevent progression of gingivitis in well-controlled diabetes.
| Biology and health sciences | Specific diseases | Health |
19644137 | https://en.wikipedia.org/wiki/Mobile%20phone | Mobile phone | A mobile phone, or cell phone, is a portable telephone that allows users to make and receive calls over a radio frequency link while moving within a designated telephone service area, unlike fixed-location phones (landline phones). This radio frequency link connects to the switching systems of a mobile phone operator, providing access to the public switched telephone network (PSTN). Modern mobile telephony relies on a cellular network architecture, which is why mobile phones are often referred to as 'cell phones' in North America.
Beyond traditional voice communication, digital mobile phones have evolved to support a wide range of additional services. These include text messaging, multimedia messaging, email, and internet access (via LTE, 5G NR or Wi-Fi), as well as short-range wireless technologies like Bluetooth, infrared, and ultra-wideband (UWB).
Mobile phones also support a variety of multimedia capabilities, such as digital photography, video recording, and gaming. In addition, they enable multimedia playback and streaming, including video content, as well as radio and television streaming. Furthermore, mobile phones offer satellite-based services, such as navigation and messaging, as well as business applications and payment solutions (via near-field communication (NFC)).
Mobile phones offering only basic features are often referred to as feature phones (slang: "dumbphones"), while those with advanced computing power are known as smartphones.
The first handheld mobile phone was demonstrated by Martin Cooper of Motorola in New York City on 3 April 1973, using a handset weighing c. 2 kilograms (4.4 lbs). In 1979, Nippon Telegraph and Telephone (NTT) launched the world's first cellular network in Japan. In 1983, the DynaTAC 8000x was the first commercially available handheld mobile phone. From 1983 to 2014, worldwide mobile phone subscriptions grew to over seven billion; enough to provide one for every person on Earth. In the first quarter of 2016, the top smartphone developers worldwide were Samsung, Apple and Huawei; smartphone sales represented 78 percent of total mobile phone sales. For feature phones , the top-selling brands were Samsung, Nokia and Alcatel.
Mobile phones are considered an important human invention as they have been one of the most widely used and sold pieces of consumer technology. The growth in popularity has been rapid in some places, for example, in the UK, the total number of mobile phones overtook the number of houses in 1999. Today, mobile phones are globally ubiquitous, and in almost half the world's countries, over 90% of the population owns at least one.
History
A handheld mobile radio telephone service was envisioned in the early stages of radio engineering. In 1917, Finnish inventor Eric Tigerstedt filed a patent for a "pocket-size folding telephone with a very thin carbon microphone". Early predecessors of cellular phones included analog radio communications from ships and trains. The race to create truly portable telephone devices began after World War II, with developments taking place in many countries. The advances in mobile telephony have been traced in successive "generations", starting with the early zeroth-generation (0G) services, such as Bell System's Mobile Telephone Service and its successor, the Improved Mobile Telephone Service. These 0G systems were not cellular, supported a few simultaneous calls, and were very expensive.
The first handheld cellular mobile phone was demonstrated by John F. Mitchell and Martin Cooper of Motorola in 1973, using a handset weighing . The first commercial automated cellular network (1G) analog was launched in Japan by Nippon Telegraph and Telephone in 1979. This was followed in 1981 by the simultaneous launch of the Nordic Mobile Telephone (NMT) system in Denmark, Finland, Norway, and Sweden. Several other countries then followed in the early to mid-1980s. These first-generation (1G) systems could support far more simultaneous calls but still used analog cellular technology. In 1983, the DynaTAC 8000x was the first commercially available handheld mobile phone.
In 1991, the second-generation (2G) digital cellular technology was launched in Finland by Radiolinja on the GSM standard. This sparked competition in the sector as the new operators challenged the incumbent 1G network operators. The GSM standard is a European initiative expressed at the CEPT ("Conférence Européenne des Postes et Telecommunications", European Postal and Telecommunications conference). The Franco-German R&D cooperation demonstrated the technical feasibility, and in 1987, a Memorandum of Understanding was signed between 13 European countries that agreed to launch a commercial service by 1991. The first version of the GSM standard had 6,000 pages. The IEEE and RSE awarded Thomas Haug and Philippe Dupuis the 2018 James Clerk Maxwell medal for their contributions to the first digital mobile telephone standard. In 2018, the GSM was used by over 5 billion people in over 220 countries. The GSM (2G) has evolved into 3G, 4G and 5G. The standardization body for GSM started at the CEPT Working Group GSM (Group Special Mobile) in 1982 under the umbrella of CEPT. In 1988, ETSI was established, and all CEPT standardization activities were transferred to ETSI. Working Group GSM became Technical Committee GSM. In 1991, it became Technical Committee SMG (Special Mobile Group) when ETSI tasked the committee with UMTS (3G). In addition to transmitting voice over digital signals, the 2G network introduced data services for mobile, starting with SMS text messages, then expanding to Multimedia Messaging Service (MMS), and mobile internet with a theoretical maximum transfer speed of 384 kbit/s (48 kB/s).
In 2001, the third-generation (3G) was launched in Japan by NTT DoCoMo on the WCDMA standard. This was followed by 3.5G or 3G+ enhancements based on the high-speed packet access (HSPA) family, allowing UMTS networks to have higher data transfer speeds and capacity. 3G is able to provide mobile broadband access of several Mbit/s to smartphones and mobile modems in laptop computers. This ensures it can be applied to mobile Internet access, VoIP, video calls, and sending large e-mail messages, as well as watching videos, typically in standard-definition quality.
By 2009, it had become clear that, at some point, 3G networks would be overwhelmed by the growth of bandwidth-intensive applications, such as streaming media. Consequently, the industry began looking to data-optimized fourth-generation (4G) technologies, with the promise of speed improvements up to tenfold over existing 3G technologies. The first publicly available LTE service was launched in Scandinavia by TeliaSonera in 2009. In the 2010s, 4G technology has found diverse applications across various sectors, showcasing its versatility in delivering high-speed wireless communication, such as mobile broadband, the internet of things (IoT), fixed wireless access, and multimedia streaming (including music, video, radio, and television).
Deployment of fifth-generation (5G) cellular networks commenced worldwide in 2019. The term "5G" was originally used in research papers and projects to denote the next major phase in mobile telecommunication standards beyond the 4G/IMT-Advanced standards. The 3GPP defines 5G as any system that adheres to the 5G NR (5G New Radio) standard. 5G can be implemented in low-band, mid-band or high-band millimeter-wave, with download speeds that can achieve gigabit-per-second (Gbit/s) range, aiming for a network latency of 1 ms. This near-real-time responsiveness and improved overall data performance are crucial for applications like online gaming, augmented and virtual reality, autonomous vehicles, IoT, and critical communication services.
Types
Smartphone
Smartphones are defined by their advanced computing capabilities, which include internet connectivity and access to a wide range of applications. The International Telecommunication Union measures those with Internet connection, which it calls Active Mobile-Broadband subscriptions (which includes tablets, etc.). In developed countries, smartphones have largely replaced earlier mobile technologies, while in developing regions, they account for around 50% of all mobile phone usage.
Feature phone
Feature phone is a term typically used as a retronym to describe mobile phones which are limited in capabilities in contrast to a modern smartphone. Feature phones typically provide voice calling and text messaging functionality, in addition to basic multimedia and Internet capabilities, and other services offered by the user's wireless service provider. A feature phone has additional functions over and above a basic mobile phone, which is only capable of voice calling and text messaging. Feature phones and basic mobile phones tend to use a proprietary, custom-designed software and user interface. By contrast, smartphones generally use a mobile operating system that often shares common traits across devices.
Infrastructure
The critical advantage that modern cellular networks have over predecessor systems is the concept of frequency reuse allowing many simultaneous telephone conversations in a given service area. This allows efficient use of the limited radio spectrum allocated to mobile services, and lets thousands of subscribers converse at the same time within a given geographic area.
Former systems would cover a service area with one or two powerful base stations with a range of up to tens of kilometers' (miles), using only a few sets of radio channels (frequencies). Once these few channels were in use by customers, no further customers could be served until another user vacated a channel. It would be impractical to give every customer a unique channel since there would not be enough bandwidth allocated to the mobile service. As well, technical limitations such as antenna efficiency and receiver design limit the range of frequencies a customer unit could use.
A cellular network mobile phone system gets its name from dividing the service area into many small cells, each with a base station with (for example) a useful range on the order of a kilometer (mile). These systems have dozens or hundreds of possible channels allocated to them. When a subscriber is using a given channel for a telephone connection, that frequency is unavailable for other customers in the local cell and in the adjacent cells. However, cells further away can re-use that channel without interference as the subscriber's handset is too far away to be detected. The transmitter power of each base station is coordinated to efficiently service its own cell, but not to interfere with the cells further away.
Automation embedded in the customer's handset and in the base stations control all phases of the call, from detecting the presence of a handset in a service area, temporary assignment of a channel to a handset making a call, interface with the land-line side of the network to connect to other subscribers, and collection of billing information for the service. The automation systems can control the "hand off" of a customer handset moving between one cell and another so that a call in progress continues without interruption, changing channels if required. In the earliest mobile phone systems by contrast, all control was done manually; the customer would search for an unoccupied channel and speak to a mobile operator to request connection of a call to a landline number or another mobile. At the termination of the call the mobile operator would manually record the billing information.
Mobile phones communicate with cell towers that are placed to give coverage across a telephone service area, which is divided up into 'cells'. Each cell uses a different set of frequencies from neighboring cells, and will typically be covered by three towers placed at different locations. The cell towers are usually interconnected to each other and the phone network and the internet by wired connections. Due to bandwidth limitations each cell will have a maximum number of cell phones it can handle at once. The cells are therefore sized depending on the expected usage density, and may be much smaller in cities. In that case much lower transmitter powers are used to avoid broadcasting beyond the cell.
In order to handle the high traffic, multiple towers can be set up in the same area (using different frequencies). This can be done permanently or temporarily such as at special events or in disasters. Cell phone companies will bring a truck with equipment to host the abnormally high traffic.
Capacity was further increased when phone companies implemented digital networks. With digital, one frequency can host multiple simultaneous calls.
Additionally, short-range Wi-Fi infrastructure is often used by smartphones as much as possible as it offloads traffic from cell networks on to local area networks.
Hardware
The common components found on all mobile phones are:
A central processing unit (CPU), the processor of phones. The CPU is a microprocessor fabricated on a metal–oxide–semiconductor (MOS) integrated circuit (IC) chip.
A battery, providing the power source for the phone functions. A modern handset typically uses a lithium-ion battery (LIB), whereas older handsets used nickel–metal hydride (Ni–MH) batteries.
An input mechanism to allow the user to interact with the phone. These are a keypad for feature phones, and touch screens for most smartphones (typically with capacitive sensing).
A display which echoes the user's typing, and displays text messages, contacts, and more. The display is typically either a liquid-crystal display (LCD) or organic light-emitting diode (OLED) display.
Speakers for sound.
Subscriber identity module (SIM) cards and removable user identity module (R-UIM) cards.
A hardware notification LED on some phones
Low-end mobile phones are often referred to as feature phones and offer basic telephony. Handsets with more advanced computing ability through the use of native software applications are known as smartphones. The first GSM phones and many feature phones had NOR flash memory, from which processor instructions could be executed directly in an execute in place architecture and allowed for short boot times. With smartphones, NAND flash memory was adopted as it has larger storage capacities and lower costs, but causes longer boot times because instructions cannot be executed from it directly, and must be copied to RAM memory first before execution.
Central processing unit
Mobile phones have central processing units (CPUs), similar to those in computers, but optimised to operate in low power environments.
Mobile CPU performance depends not only on the clock rate (generally given in multiples of hertz) but also the memory hierarchy also greatly affects overall performance. Because of these problems, the performance of mobile phone CPUs is often more appropriately given by scores derived from various standardized tests to measure the real effective performance in commonly used applications.
Display
One of the main characteristics of phones is the screen. Depending on the device's type and design, the screen fills most or nearly all of the space on a device's front surface. Many smartphone displays have an aspect ratio of 16:9, but taller aspect ratios became more common in 2017.
Screen sizes are often measured in diagonal inches or millimeters; feature phones generally have screen sizes below . Phones with screens larger than are often called "phablets." Smartphones with screens over in size are commonly difficult to use with only a single hand, since most thumbs cannot reach the entire screen surface; they may need to be shifted around in the hand, held in one hand and manipulated by the other, or used in place with both hands. Due to design advances, some modern smartphones with large screen sizes and "edge-to-edge" designs have compact builds that improve their ergonomics, while the shift to taller aspect ratios have resulted in phones that have larger screen sizes whilst maintaining the ergonomics associated with smaller 16:9 displays.
Liquid-crystal displays are the most common; others are IPS, LED, OLED, and AMOLED displays. Some displays are integrated with pressure-sensitive digitizers, such as those developed by Wacom and Samsung, and Apple's "3D Touch" system.
Sound
In sound, smartphones and feature phones vary little. Some audio-quality enhancing features, such as Voice over LTE and HD Voice, have appeared and are often available on newer smartphones. Sound quality can remain a problem due to the design of the phone, the quality of the cellular network and compression algorithms used in long-distance calls. Audio quality can be improved using a VoIP application over WiFi. Cellphones have small speakers so that the user can use a speakerphone feature and talk to a person on the phone without holding it to their ear. The small speakers can also be used to listen to digital audio files of music or speech or watch videos with an audio component, without holding the phone close to the ear.
Battery
The typical lifespan of a mobile phone battery is approximately two to three years, although this varies based on usage patterns, environmental conditions, and overall care. Most modern mobile phones use lithium-ion (Li-ion) batteries, which are designed to endure between 500 and 2,500 charge cycles. The exact number of cycles depends on factors such as charging habits, operating temperature, and battery management systems.
Li-ion batteries gradually degrade over time due to chemical aging, leading to reduced capacity and performance, often noticeable after one or two years of regular use. Unlike older battery types, such as nickel-metal hydride (Ni-MH), Li-ion batteries do not need to be fully discharged to maintain their longevity. In fact, they perform best when kept between 30% and 80% of their full charge. While practices such as avoiding excessive heat and minimizing overcharging can help preserve battery health, many modern devices include built-in safeguards. These safeguards, typically managed by the phone's internal battery management system (BMS), prevent overcharging by cutting off power once the battery reaches full capacity. Additionally, most contemporary chargers and devices are designed to regulate charging to minimize stress on the battery. Therefore, while good charging habits can positively impact battery longevity, most users benefit from these integrated protections, making battery maintenance less of a concern in day-to-day use.
Future mobile phone batteries are expected to utilize advanced technologies such as silicon-carbon (Si/C) batteries and solid-state batteries, which promise to offer higher energy densities, longer lifespans, and improved safety compared to current lithium-ion batteries.
SIM card
Mobile phones require a small microchip called a Subscriber Identity Module or SIM card, in order to function. The SIM card is approximately the size of a small postage stamp and is usually placed underneath the battery in the rear of the unit. The SIM securely stores the service-subscriber key (IMSI) and the Ki used to identify and authenticate the user of the mobile phone. The SIM card allows users to change phones by simply removing the SIM card from one mobile phone and inserting it into another mobile phone or broadband telephony device, provided that this is not prevented by a SIM lock. The first SIM card was made in 1991 by Munich smart card maker Giesecke & Devrient for the Finnish wireless network operator Radiolinja.
A hybrid mobile phone can hold up to four SIM cards, with a phone having a different device identifier for each SIM Card. SIM and R-UIM cards may be mixed together to allow both GSM and CDMA networks to be accessed. From 2010 onwards, such phones became popular in emerging markets, and this was attributed to the desire to obtain the lowest calling costs.
When the removal of a SIM card is detected by the operating system, it may deny further operation until a reboot.
Software
Software platforms
Feature phones have basic software platforms. Smartphones have advanced software platforms. Android OS has been the best-selling OS worldwide on smartphones since 2011.
Mobile app
A mobile app is a computer program designed to run on a mobile device, such as a smartphone. The term "app" is a shortening of the term "software application".
Messaging
A common data application on mobile phones is Short Message Service (SMS) text messaging. The first SMS message was sent from a computer to a mobile phone in 1992 in the UK while the first person-to-person SMS from phone to phone was sent in Finland in 1993. The first mobile news service, delivered via SMS, was launched in Finland in 2000, and subsequently many organizations provided "on-demand" and "instant" news services by SMS. Multimedia Messaging Service (MMS) was introduced in March 2002.
Application stores
The introduction of Apple's App Store for the iPhone and iPod Touch in July 2008 popularized manufacturer-hosted online distribution for third-party applications (software and computer programs) focused on a single platform. There are a huge variety of apps, including video games, music products and business tools. Up until that point, smartphone application distribution depended on third-party sources providing applications for multiple platforms, such as GetJar, Handango, Handmark, and PocketGear. Following the success of the App Store, other smartphone manufacturers launched application stores, such as Google's Android Market (later renamed to the Google Play Store), RIM's BlackBerry App World, or Android-related app stores like Aptoide, Cafe Bazaar, F-Droid, GetJar, and Opera Mobile Store. In February 2014, 93% of mobile developers were targeting smartphones first for mobile app development.
Sales
By manufacturer
As of 2022, the top five manufacturers worldwide were Samsung (21%), Apple (16%), Xiaomi (13%), Oppo (10%), and Vivo (9%).
History
From 1983 to 1998, Motorola was market leader in mobile phones. Nokia was the market leader in mobile phones from 1998 to 2012. In Q1 2012, Samsung surpassed Nokia, selling 93.5 million units as against Nokia's 82.7 million units. Samsung has retained its top position since then.
Aside from Motorola, European brands such as Nokia, Siemens and Ericsson once held large sway over the global mobile phone market, and many new technologies were pioneered in Europe. By 2010, the influence of European companies had significantly decreased due to fierce competition from American and Asian companies, to where most technical innovation had shifted. Apple and Google, both of the United States, also came to dominate mobile phone software.
By mobile phone operator
The world's largest individual mobile operator by number of subscribers is China Mobile, which has over 902 million mobile phone subscribers . Over 50 mobile operators have over ten million subscribers each, and over 150 mobile operators had at least one million subscribers by the end of 2009. In 2014, there were more than seven billion mobile phone subscribers worldwide, a number that is expected to keep growing.
Use
Mobile phones are used for a variety of purposes, such as keeping in touch with family members, for conducting business, and in order to have access to a telephone in the event of an emergency. Some people carry more than one mobile phone for different purposes, such as for business and personal use. Multiple SIM cards may be used to take advantage of the benefits of different calling plans. For example, a particular plan might provide for cheaper local calls, long-distance calls, international calls, or roaming.
The mobile phone has been used in a variety of diverse contexts in society. For example:
A study by Motorola found that one in ten mobile phone subscribers have a second phone that is often kept secret from other family members. These phones may be used to engage in such activities as extramarital affairs or clandestine business dealings.
Some organizations assist victims of domestic violence by providing mobile phones for use in emergencies. These are often refurbished phones.
The advent of widespread text-messaging has resulted in the cell phone novel, the first literary genre to emerge from the cellular age, via text messaging to a website that collects the novels as a whole.
Mobile telephony also facilitates activism and citizen journalism.
The United Nations reported that mobile phones have spread faster than any other form of technology and can improve the livelihood of the poorest people in developing countries, by providing access to information in places where landlines or the Internet are not available, especially in the least developed countries. Use of mobile phones also spawns a wealth of micro-enterprises, by providing such work as selling airtime on the streets and repairing or refurbishing handsets.
In Mali and other African countries, people used to travel from village to village to let friends and relatives know about weddings, births, and other events. This can now be avoided in areas with mobile phone coverage, which are usually more extensive than areas with just land-line penetration.
The TV industry has recently started using mobile phones to drive live TV viewing through mobile apps, advertising, social TV, and mobile TV. It is estimated that 86% of Americans use their mobile phone while watching TV.
In some parts of the world, mobile phone sharing is common. Cell phone sharing is prevalent in urban India, as families and groups of friends often share one or more mobile phones among their members. There are obvious economic benefits, but often familial customs and traditional gender roles play a part. It is common for a village to have access to only one mobile phone, perhaps owned by a teacher or missionary, which is available to all members of the village for necessary calls.
Smartphones also have the use for individuals who suffer from diabetes. There are apps for patients with diabetes to self monitor their blood sugar, and can sync with flash monitors. The apps have a feature to send automated feedback or possible warnings to other family members or healthcare providers in the case of an emergency.
Content distribution
In 1998, one of the first examples of distributing and selling media content through the mobile phone was the sale of ringtones by Radiolinja in Finland. Soon afterwards, other media content appeared, such as news, video games, jokes, horoscopes, TV content and advertising. Most early content for mobile phones tended to be copies of legacy media, such as banner advertisements or TV news highlight video clips. Recently, unique content for mobile phones has been emerging, from ringtones and ringback tones to mobisodes, video content that has been produced exclusively for mobile phones.
Mobile banking and payment
In many countries, mobile phones are used to provide mobile banking services, which may include the ability to transfer cash payments by secure SMS text message. Kenya's M-PESA mobile banking service, for example, allows customers of the mobile phone operator Safaricom to hold cash balances which are recorded on their SIM cards. Cash can be deposited or withdrawn from M-PESA accounts at Safaricom retail outlets located throughout the country and can be transferred electronically from person to person and used to pay bills to companies.
Branchless banking has also been successful in South Africa and the Philippines. A pilot project in Bali was launched in 2011 by the International Finance Corporation and an Indonesian bank, Bank Mandiri.
Mobile payments were first trialled in Finland in 1998 when two Coca-Cola vending machines in Espoo were enabled to work with SMS payments. Eventually, the idea spread and in 1999, the Philippines launched the country's first commercial mobile payments systems with mobile operators Globe and Smart.
Some mobile phones can make mobile payments via direct mobile billing schemes, or through contactless payments if the phone and the point of sale support near field communication (NFC). Enabling contactless payments through NFC-equipped mobile phones requires the co-operation of manufacturers, network operators, and retail merchants.
Mobile tracking
Mobile phones are commonly used to collect location data. While the phone is turned on, the geographical location of a mobile phone can be determined easily (whether it is being used or not) using a technique known as multilateration to calculate the differences in time for a signal to travel from the mobile phone to each of several cell towers near the owner of the phone.
The movements of a mobile phone user can be tracked by their service provider and, if desired, by law enforcement agencies and their governments. Both the SIM card and the handset can be tracked.
China has proposed using this technology to track the commuting patterns of Beijing city residents. In the UK and US, law enforcement and intelligence services use mobile phones to perform surveillance operations.
Hackers have been able to track a phone's location, read messages, and record calls, through obtaining a subscribers phone number.
Electronic waste regulation
Studies have shown that around 40–50% of the environmental impact of mobile phones occurs during the manufacture of their printed wiring boards and integrated circuits.
The average user replaces their mobile phone every 11 to 18 months, and the discarded phones then contribute to electronic waste. Mobile phone manufacturers within Europe are subject to the WEEE directive, and Australia has introduced a mobile phone recycling scheme.
Apple Inc. had an advanced robotic disassembler and sorter called Liam specifically for recycling outdated or broken iPhones.
Theft
According to the Federal Communications Commission, one out of three robberies involve the theft of a cellular phone. Police data in San Francisco show that half of all robberies in 2012 were thefts of cellular phones. An online petition on Change.org, called Secure our Smartphones, urged smartphone manufacturers to install kill switches in their devices to make them unusable if stolen. The petition is part of a joint effort by New York Attorney General Eric Schneiderman and San Francisco District Attorney George Gascón and was directed to the CEOs of the major smartphone manufacturers and telecommunication carriers. On 10 June 2013, Apple announced that it would install a "kill switch" on its next iPhone operating system, due to debut in October 2013.
All mobile phones have a unique identifier called IMEI. Anyone can report their phone as lost or stolen with their Telecom Carrier, and the IMEI would be blacklisted with a central registry. Telecom carriers, depending upon local regulation can or must implement blocking of blacklisted phones in their network. There are, however, a number of ways to circumvent a blacklist. One method is to send the phone to a country where the telecom carriers are not required to implement the blacklisting and sell it there, another involves altering the phone's IMEI number. Even so, mobile phones typically have less value on the second-hand market if the phones original IMEI is blacklisted.
Conflict minerals
Demand for metals used in mobile phones and other electronics fuelled the Second Congo War, which claimed almost 5.5 million lives. In a 2012 news story, The Guardian reported: "In unsafe mines deep underground in eastern Congo, children are working to extract minerals essential for the electronics industry. The profits from the minerals finance the bloodiest conflict since the second world war; the war has lasted nearly 20 years and has recently flared up again. For the last 15 years, the Democratic Republic of the Congo has been a major source of natural resources for the mobile phone industry." The company Fairphone has worked to develop a mobile phone that does not contain conflict minerals.
Kosher phones
Due to concerns by the Orthodox Jewish rabbinate in Britain that texting by youths could waste time and lead to "immodest" communication, the rabbinate recommended that phones with text-messaging capability not be used by children; to address this, they gave their official approval to a brand of "Kosher" phones with no texting capabilities. Although these phones are intended to prevent immodesty, some vendors report good sales to adults who prefer the simplicity of the devices; other Orthodox Jews question the need for them.
In Israel, similar phones to kosher phones with restricted features exist to observe the sabbath; under Orthodox Judaism, the use of any electrical device is generally prohibited during this time, other than to save lives, or reduce the risk of death or similar needs. Such phones are approved for use by essential workers, such as health, security, and public service workers.
Restrictions
Restrictions on the use of mobile phones are applied in a number of different contexts, often with the goal of health, safety, security or proper functioning of an establishment, or as a matter of etiquette. Such contexts include:
While driving
Mobile phone use while driving, including talking on the phone, texting, or operating other phone features, is common but controversial. It is widely considered dangerous due to distracted driving. Being distracted while operating a motor vehicle has been shown to increase the risk of accidents. In September 2010, the US National Highway Traffic Safety Administration (NHTSA) reported that 995 people were killed by drivers distracted by cell phones. In March 2011, a US insurance company, State Farm Insurance, announced the results of a study which showed 19% of drivers surveyed accessed the Internet on a smartphone while driving. Many jurisdictions prohibit the use of mobile phones while driving. In Egypt, Israel, Japan, Portugal, and Singapore, both handheld and hands-free use of a mobile phone (which uses a speakerphone) is banned. In other countries, including the UK and France and in many US states, only handheld phone use is banned while hands-free use is permitted.
A 2011 study reported that over 90% of college students surveyed text (initiate, reply or read) while driving.
The scientific literature on the dangers of driving while sending a text message from a mobile phone, or texting while driving, is limited. A simulation study at the University of Utah found a sixfold increase in distraction-related accidents when texting.
Due to the increasing complexity of mobile phones, they are often more like mobile computers in their available uses. This has introduced additional difficulties for law enforcement officials when attempting to distinguish one usage from another in drivers using their devices. This is more apparent in countries which ban both handheld and hands-free usage, rather than those which ban handheld use only, as officials cannot easily tell which function of the mobile phone is being used simply by looking at the driver. This can lead to drivers being stopped for using their device illegally for a phone call when, in fact, they were using the device legally, for example, when using the phone's incorporated controls for car stereo, GPS or satnav.
A 2010 study reviewed the incidence of mobile phone use while cycling and its effects on behaviour and safety. In 2013, a national survey in the US reported the number of drivers who reported using their cellphones to access the Internet while driving had risen to nearly one of four. A study conducted by the University of Vienna examined approaches for reducing inappropriate and problematic use of mobile phones, such as using mobile phones while driving.
Accidents involving a driver being distracted by talking on a mobile phone have begun to be prosecuted as negligence similar to speeding. In the United Kingdom, from 27 February 2007, motorists who are caught using a hand-held mobile phone while driving will have three penalty points added to their license in addition to the fine of £60. This increase was introduced to try to stem the increase in drivers ignoring the law. Japan prohibits all mobile phone use while driving, including use of hands-free devices. New Zealand has banned hand-held cell phone use since 1 November 2009. Many states in the United States have banned texting on cell phones while driving. Illinois became the 17th American state to enforce this law. , 30 states had banned texting while driving, with Kentucky becoming the most recent addition on 15 July.
Public Health Law Research maintains a list of distracted driving laws in the United States. This database of laws provides a comprehensive view of the provisions of laws that restrict the use of mobile communication devices while driving for all 50 states and the District of Columbia between 1992 when first law was passed, through 1 December 2010. The dataset contains information on 22 dichotomous, continuous or categorical variables including, for example, activities regulated (e.g., texting versus talking, hands-free versus handheld), targeted populations, and exemptions.
On aircraft
While walking
In 2010, an estimated 1500 pedestrians were injured in the US while using a cellphone and some jurisdictions have attempted to ban pedestrians from using their cellphones. Other countries, such as China and the Netherlands, have introduced special lanes for smartphone users to help direct and manage them.
In prisons
In hospitals
As of 2007, some hospitals had banned mobile devices due to a common misconception that their use would create significant electromagnetic interference.
Health effects
Screen time, the amount of time using a device with a screen, has become an issue for mobile phones since the adaptation of smartphones. Research is being conducted to show the correlation between screen time and the mental and physical harm in child development. To prevent harm, some parents and even governments have placed restrictions on its usage.
There have been rumors that mobile phone use can cause cancer, but this is a myth.
While there are rumors of mobile phones causing cancer, there was a study conducted by International Agency for Research on Cancer (IARC) that stated the there could be an increase risk of brain tumors with the use of smartphones, this is not confirmed. They also stated that with the lack of data for the research and the usage periods of 15 years will warrant further research for smartphones and the cause of brain tumors.
Educational impact
A study by the London School of Economics found that banning mobile phones in schools could increase pupils' academic performance, providing benefits equal to one extra week of schooling per year.
Culture and popularity
Mobile phones are considered an important human invention as it has been one of the most widely used and sold pieces of consumer technology. They have also become culturally symbolic. In Japanese mobile phone culture for example, mobile phones are often decorated with charms. They have also become fashion symbols at times. The Motorola Razr V3 and LG Chocolate are two examples of devices that were popular for being fashionable while not necessarily focusing on the original purpose of mobile phones, i.e. a device to provide mobile telephony.
Some have also suggested that mobile phones or smartphones are a status symbol. For example a research paper suggested that owning specifically an Apple iPhone was seen to be a status symbol.
Text messaging, which are performed on mobile phones, has also led to the creation of 'SMS language'. It also led to the growing popularity of emojis.
| Technology | Media and communication | null |
2589714 | https://en.wikipedia.org/wiki/Milky%20Way | Milky Way | The Milky Way is the galaxy that includes the Solar System, with the name describing the galaxy's appearance from Earth: a hazy band of light seen in the night sky formed from stars that cannot be individually distinguished by the naked eye.
The Milky Way is a barred spiral galaxy with a D25 isophotal diameter estimated at , but only about 1,000 light-years thick at the spiral arms (more at the bulge). Recent simulations suggest that a dark matter area, also containing some visible stars, may extend up to a diameter of almost 2 million light-years (613 kpc). The Milky Way has several satellite galaxies and is part of the Local Group of galaxies, forming part of the Virgo Supercluster which is itself a component of the Laniakea Supercluster.
It is estimated to contain 100–400 billion stars and at least that number of planets. The Solar System is located at a radius of about 27,000 light-years (8.3 kpc) from the Galactic Center, on the inner edge of the Orion Arm, one of the spiral-shaped concentrations of gas and dust. The stars in the innermost 10,000 light-years form a bulge and one or more bars that radiate from the bulge. The Galactic Center is an intense radio source known as Sagittarius A*, a supermassive black hole of 4.100 (± 0.034) million solar masses. The oldest stars in the Milky Way are nearly as old as the Universe itself and thus probably formed shortly after the Dark Ages of the Big Bang.
Galileo Galilei first resolved the band of light into individual stars with his telescope in 1610. Until the early 1920s, most astronomers thought that the Milky Way contained all the stars in the Universe. Following the 1920 Great Debate between the astronomers Harlow Shapley and Heber Doust Curtis, observations by Edwin Hubble in 1923 showed that the Milky Way is just one of many galaxies.
Etymology and mythology
In the Babylonian epic poem Enūma Eliš, the Milky Way is created from the severed tail of the primeval salt water dragoness Tiamat, set in the sky by Marduk, the Babylonian national god, after slaying her. This story was once thought to have been based on an older Sumerian version in which Tiamat is instead slain by Enlil of Nippur, but is now thought to be purely an invention of Babylonian propagandists with the intention to show Marduk as superior to the Sumerian deities.
In Greek mythology, Zeus places his son born by a mortal woman, the infant Heracles, on Hera's breast while she is asleep so the baby will drink her divine milk and become immortal. Hera wakes up while breastfeeding and then realizes she is nursing an unknown baby: she pushes the baby away, some of her milk spills, and it produces the band of light known as the Milky Way. In another Greek story, the abandoned Heracles is given by Athena to Hera for feeding, but Heracles' forcefulness causes Hera to rip him from her breast in pain.
Llys Dôn (literally "The Court of Dôn") is the traditional Welsh name for the constellation Cassiopeia. At least two of Dôn's children also have astronomical associations: Caer Gwydion ("The fortress of Gwydion") is the traditional Welsh name for the Milky Way, and Caer Arianrhod ("The Fortress of Arianrhod") being the constellation of Corona Borealis.
In Western culture, the name "Milky Way" is derived from its appearance as a dim un-resolved "milky" glowing band arching across the night sky. The term is a translation of the Classical Latin via lactea, in turn derived from the Hellenistic Greek , short for (), meaning "milky circle". The Ancient Greek () – from root -, ("milk") + (forming adjectives) – is also the root of "galaxy", the name for our, and later all such, collections of stars.
The Milky Way, or "milk circle", was just one of 11 "circles" the Greeks identified in the sky, others being the zodiac, the meridian, the horizon, the equator, the tropics of Cancer and Capricorn, the Arctic Circle and the Antarctic Circle, and two colure circles passing through both poles.
Common names
"Birds' Path" is used in several Uralic and Turkic languages and in the Baltic languages. Northern peoples observed that migratory birds follow the course of the galaxy while migrating at the Northern Hemisphere. The name "Birds' Path" (in Finnish, Estonian, Latvian, Lithuanian, Bashkir and Kazakh) has some variations in other languages, e.g. "Way of the grey (wild) goose" in Chuvash, Mari and Tatar and "Way of the Crane" in Erzya and Moksha.
House river: The Kaurna people of the Adelaide Plains of South Australia called the Milky Way wodliparri in the Kaurna language, meaning "house river".
Emu in the Sky: The Gomeroi people between New South Wales and Queensland called the Milky Way Dhinawan, the giant Emu in the Sky that it stretches across the night sky.
Milky Way: Many European languages have borrowed, directly or indirectly, the Greek name for the Milky Way, including English and Latin.
Road to Santiago: the Milky Way was traditionally used as a guide by pilgrims traveling to the holy site at Compostela, hence the use of "The Road to Santiago" as a name for the Milky Way. Curiously, La Voje Ladee "The Milky Way" was also used to refer to the pilgrimage road.
River Ganga of the Sky: this Sanskrit name ( Ākāśagaṃgā) is used in many Indian languages following a Hindu belief .
Silver River: this Chinese name "Silver River" () is used throughout East Asia, including Korea and Vietnam. In Japan and Korea, "Silver River" (; ) means galaxies in general.
River of Heaven: The Japanese name for the Milky Way is the "River of Heaven" (, Amanokawa), as well as an alternative name in Chinese ().
Straw Way: In West Asia, Central Asia and parts of the Balkans the name for the Milky Way is related to the word for straw. Today, Persians, Pakistanis, and Turks use it in addition to Arabs. It has been suggested that the term was spread by medieval Arabs who in turn borrowed it from Armenians.
Walsingham Way: In England the Milky Way was called the Walsingham Way in reference to the shrine of Our Lady of Walsingham which is in Norfolk, England. It was understood to be either a guide to the pilgrims who flocked there, or a representation of the pilgrims themselves.
Winter Street: Scandinavian peoples, such as Swedes, have called the galaxy Winter Street (Vintergatan) as the galaxy is most clearly visible during the winter at the northern hemisphere, especially at high latitudes where the glow of the Sun late at night can obscure it during the summer.
Appearance
The Milky Way is visible as a hazy band of white light, some 30° wide, arching the night sky. Although all the individual naked-eye stars in the entire sky are part of the Milky Way Galaxy, the term "Milky Way" is limited to this band of light. The light originates from the accumulation of unresolved stars and other material located in the direction of the galactic plane. Brighter regions around the band appear as soft visual patches known as star clouds. The most conspicuous of these is the Large Sagittarius Star Cloud, a portion of the central bulge of the galaxy. Dark regions within the band, such as the Great Rift and the Coalsack, are areas where interstellar dust blocks light from distant stars. Peoples of the southern hemisphere, including the Inca and Australian aborigines, identified these regions as dark cloud constellations. The area of sky that the Milky Way obscures is called the Zone of Avoidance.
The Milky Way has a relatively low surface brightness. Its visibility can be greatly reduced by background light, such as light pollution or moonlight. The sky needs to be darker than about 20.2 magnitude per square arcsecond in order for the Milky Way to be visible. It should be visible if the limiting magnitude is approximately +5.1 or better and shows a great deal of detail at +6.1. This makes the Milky Way difficult to see from brightly lit urban or suburban areas, but very prominent when viewed from rural areas when the Moon is below the horizon. Maps of artificial night sky brightness show that more than one-third of Earth's population cannot see the Milky Way from their homes due to light pollution.
As viewed from Earth, the visible region of the Milky Way's galactic plane occupies an area of the sky that includes 30 constellations. The Galactic Center lies in the direction of Sagittarius, where the Milky Way is brightest. From Sagittarius, the hazy band of white light appears to pass around to the galactic anticenter in Auriga. The band then continues the rest of the way around the sky, back to Sagittarius, dividing the sky into two roughly equal hemispheres.
The galactic plane is inclined by about 60° to the ecliptic (the plane of Earth's orbit). Relative to the celestial equator, it passes as far north as the constellation of Cassiopeia and as far south as the constellation of Crux, indicating the high inclination of Earth's equatorial plane and the plane of the ecliptic, relative to the galactic plane. The north galactic pole is situated at right ascension 12h 49m, declination +27.4° (B1950) near β Comae Berenices, and the south galactic pole is near α Sculptoris. Because of this high inclination, depending on the time of night and year, the Milky Way arch may appear relatively low or relatively high in the sky. For observers from latitudes approximately 65° north to 65° south, the Milky Way passes directly overhead twice a day.
Astronomical history
Ancient, naked eye observations
In Meteorologica, Aristotle (384–322 BC) states that the Greek philosophers Anaxagoras (–428 BC) and Democritus (460–370 BC) proposed that the Milky Way is the glow of stars not directly visible due to Earth's shadow, while other stars receive their light from the Sun, but have their glow obscured by solar rays. Aristotle himself believed that the Milky Way was part of the Earth's upper atmosphere, along with the stars, and that it was a byproduct of stars burning that did not dissipate because of its outermost location in the atmosphere, composing its great circle. He said that the milky appearance of the Milky Way Galaxy is due to the refraction of the Earth's atmosphere. The Neoplatonist philosopher Olympiodorus the Younger (–570 AD) criticized this view, arguing that if the Milky Way were sublunary, it should appear different at different times and places on Earth, and that it should have parallax, which it does not. In his view, the Milky Way is celestial. This idea would be influential later in the Muslim world.
The Persian astronomer Al-Biruni (973–1048) proposed that the Milky Way is "a collection of countless fragments of the nature of nebulous stars". The Andalusian astronomer Avempace ( 1138) proposed that the Milky Way was made up of many stars but appeared to be a continuous image in the Earth's atmosphere, citing his observation of a conjunction of Jupiter and Mars in 1106 or 1107 as evidence. The Persian astronomer Nasir al-Din al-Tusi (1201–1274) in his Tadhkira wrote: "The Milky Way, i.e. the Galaxy, is made up of a very large number of small, tightly clustered stars, which, on account of their concentration and smallness, seem to be cloudy patches. Because of this, it was likened to milk in color." Ibn Qayyim al-Jawziyya (1292–1350) proposed that the Milky Way is "a myriad of tiny stars packed together in the sphere of the fixed stars".
Telescopic observations
Proof of the Milky Way consisting of many stars came in 1610 when Galileo Galilei used a telescope to study the Milky Way and discovered that it is composed of a huge number of faint stars. Galileo also concluded that the appearance of the Milky Way was due to refraction of the Earth's atmosphere. In a treatise in 1755, Immanuel Kant, drawing on earlier work by Thomas Wright, speculated (correctly) that the Milky Way might be a rotating body of a huge number of stars, held together by gravitational forces akin to the Solar System but on much larger scales. The resulting disk of stars would be seen as a band on the sky from our perspective inside the disk. Wright and Kant also conjectured that some of the nebulae visible in the night sky might be separate "galaxies" themselves, similar to our own. Kant referred to both the Milky Way and the "extragalactic nebulae" as "island universes", a term still current up to the 1930s.
The first attempt to describe the shape of the Milky Way and the position of the Sun within it was carried out by William Herschel in 1785 by carefully counting the number of stars in different regions of the visible sky. He produced a diagram of the shape of the Milky Way with the Solar System close to the center.
In 1845, Lord Rosse constructed a new telescope and was able to distinguish between elliptical and spiral-shaped nebulae. He also managed to make out individual point sources in some of these nebulae, lending credence to Kant's earlier conjecture.
In 1904, studying the proper motions of stars, Jacobus Kapteyn reported that these were not random, as it was believed in that time; stars could be divided into two streams, moving in nearly opposite directions. It was later realized that Kapteyn's data had been the first evidence of the rotation of our galaxy, which ultimately led to the finding of galactic rotation by Bertil Lindblad and Jan Oort.
In 1917, Heber Doust Curtis had observed the nova S Andromedae within the Great Andromeda Nebula (Messier object 31). Searching the photographic record, he found 11 more novae. Curtis noticed that these novae were, on average, 10 magnitudes fainter than those that occurred within the Milky Way. As a result, he was able to come up with a distance estimate of 150,000 parsecs. He became a proponent of the "island universes" hypothesis, which held that the spiral nebulae were independent galaxies. In 1920 the Great Debate took place between Harlow Shapley and Heber Curtis, concerning the nature of the Milky Way, spiral nebulae, and the dimensions of the Universe. To support his claim that the Great Andromeda Nebula is an external galaxy, Curtis noted the appearance of dark lanes resembling the dust clouds in the Milky Way, as well as the significant Doppler shift.
The controversy was conclusively settled by Edwin Hubble in the early 1920s using the Mount Wilson observatory Hooker telescope. With the light-gathering power of this new telescope, he was able to produce astronomical photographs that resolved the outer parts of some spiral nebulae as collections of individual stars. He was also able to identify some Cepheid variables that he could use as a benchmark to estimate the distance to the nebulae. He found that the Andromeda Nebula is 275,000 parsecs from the Sun, far too distant to be part of the Milky Way.
Satellite observations
The ESA spacecraft Gaia provides distance estimates by determining the parallax of a billion stars and is mapping the Milky Way with four planned releases of maps in 2016, 2018, 2021 and 2024.
Data from Gaia has been described as "transformational". It has been estimated that Gaia has expanded the number of observations of stars from about 2 million stars as of the 1990s to 2 billion. It has expanded the measurable volume of space by a factor of 100 in radius and a factor of 1,000 in precision.
A study in 2020 concluded that Gaia detected a wobbling motion of the galaxy, which might be caused by "torques from a misalignment of the disc's rotation axis with respect to the principal axis of a non-spherical halo, or from accreted matter in the halo acquired during late infall, or from nearby, interacting satellite galaxies and their consequent tides". In April 2024, initial studies (and related maps) involving the magnetic fields of the Milky Way were reported.
Astrography
Sun's location and neighborhood
The Sun is near the inner rim of the Orion Arm, within the Local Fluff of the Local Bubble, between the Radcliffe wave and Split linear structures (formerly Gould Belt). Based upon studies of stellar orbits around Sgr A* by Gillessen et al. (2016), the Sun lies at an estimated distance of from the Galactic Center. Boehle et al. (2016) found a smaller value of , also using a star orbit analysis. The Sun is currently above, or north of, the central plane of the Galactic disk. The distance between the local arm and the next arm out, the Perseus Arm, is about . The Sun, and thus the Solar System, is located in the Milky Way's galactic habitable zone.
There are about 208 stars brighter than absolute magnitude 8.5 within a sphere with a radius of from the Sun, giving a density of one star per 69 cubic parsecs, or one star per 2,360 cubic light-years (from List of nearest bright stars). On the other hand, there are 64 known stars (of any magnitude, not counting 4 brown dwarfs) within of the Sun, giving a density of about one star per 8.2 cubic parsecs, or one per 284 cubic light-years (from List of nearest stars). This illustrates the fact that there are far more faint stars than bright stars: in the entire sky, there are about 500 stars brighter than apparent magnitude 4 but 15.5 million stars brighter than apparent magnitude 14.
The apex of the Sun's way, or the solar apex, is the direction that the Sun travels through space in the Milky Way. The general direction of the Sun's Galactic motion is towards the star Vega near the constellation of Hercules, at an angle of roughly 60 sky degrees to the direction of the Galactic Center. The Sun's orbit about the Milky Way is expected to be roughly elliptical with the addition of perturbations due to the Galactic spiral arms and non-uniform mass distributions. In addition, the Sun passes through the Galactic plane approximately 2.7 times per orbit. This is very similar to how a simple harmonic oscillator works with no drag force (damping) term. These oscillations were until recently thought to coincide with mass lifeform extinction periods on Earth. A reanalysis of the effects of the Sun's transit through the spiral structure based on CO data has failed to find a correlation.
It takes the Solar System about 240 million years to complete one orbit of the Milky Way (a galactic year), so the Sun is thought to have completed 18–20 orbits during its lifetime and 1/1250 of a revolution since the origin of humans. The orbital speed of the Solar System about the center of the Milky Way is approximately or 0.073% of the speed of light. The Sun moves through the heliosphere at . At this speed, it takes around 1,400 years for the Solar System to travel a distance of 1 light-year, or 8 days to travel 1 AU (astronomical unit). The Solar System is headed in the direction of the zodiacal constellation Scorpius, which follows the ecliptic.
Galactic quadrants
A galactic quadrant, or quadrant of the Milky Way, refers to one of four circular sectors in the division of the Milky Way. In astronomical practice, the delineation of the galactic quadrants is based upon the galactic coordinate system, which places the Sun as the origin of the mapping system.
Quadrants are described using ordinalsfor example, "1st galactic quadrant", "second galactic quadrant", or "third quadrant of the Milky Way". Viewing from the north galactic pole with 0° (zero degrees) as the ray that runs starting from the Sun and through the Galactic Center, the quadrants are:
{|
! Galacticquadrant
! Galacticlongitude(ℓ)
! Reference
|-
!align="center"|1st
|align="right" | 0° ≤ ℓ ≤ 90°
|align="center"|
|-
!align="center"|2nd
|align="right" | 90° ≤ ℓ ≤ 180°
|align="center"|
|-
!align="center"|3rd
|align="right" |180° ≤ ℓ ≤ 270°
|align="center"|
|-
!align="center"|4th
|align="right" |270° ≤ ℓ ≤ 360°(360° ≅ 0°)
|align="center"|
|}
with the galactic longitude (ℓ) increasing in the counter-clockwise direction (positive rotation) as viewed from north of the Galactic Center (a view-point several hundred thousand light-years distant from Earth in the direction of the constellation Coma Berenices); if viewed from south of the Galactic Center (a view-point similarly distant in the constellation Sculptor), ℓ would increase in the clockwise direction (negative rotation).
Size and mass
Size
The Milky Way is one of the two largest galaxies in the Local Group (the other being the Andromeda Galaxy), although the size for its galactic disc and how much it defines the isophotal diameter is not well understood. It is estimated that the significant bulk of stars in the galaxy lies within the diameter, and that the number of stars beyond the outermost disc dramatically reduces to a very low number, with respect to an extrapolation of the exponential disk with the scale length of the inner disc.
There are several methods being used in astronomy in defining the size of a galaxy, and each of them can yield different results with respect to one another. The most commonly employed method is the D25 standard – the isophote where the photometric brightness of a galaxy in the B-band (445 nm wavelength of light, in the blue part of the visible spectrum) reaches 25 mag/arcsec2. An estimate from 1997 by Goodwin and others compared the distribution of Cepheid variable stars in 17 other spiral galaxies to the ones in the Milky Way, and modelling the relationship to their surface brightnesses. This gave an isophotal diameter for the Milky Way at , by assuming that the galactic disc is well represented by an exponential disc and adopting a central surface brightness of the galaxy (μ0) of B-mag/arcsec−2 and a disk scale length (h) of .
This is significantly smaller than the Andromeda Galaxy's isophotal diameter, and slightly below the mean isophotal sizes of the galaxies being at . The paper concludes that the Milky Way and Andromeda Galaxy were not overly large spiral galaxies, nor were among the largest known (if the former not being the largest) as previously widely believed, but rather average ordinary spiral galaxies. To compare the relative physical scale of the Milky Way, if the Solar System out to Neptune were the size of a US quarter (), the Milky Way would be approximately at least the greatest north–south line of the contiguous United States. An even older study from 1978 gave a lower diameter for Milky Way about .
A 2015 paper discovered that there is a ring-like filament of stars called Triangulum–Andromeda Ring (TriAnd Ring) rippling above and below the relatively flat galactic plane, which alongside Monoceros Ring were both suggested to be primarily the result of disk oscillations and wrapping around the Milky Way, at a diameter of at least , which may be part of the Milky Way's outer disk itself, hence making the stellar disk larger by increasing to this size. A more recent 2018 paper later somewhat ruled out this hypothesis, and supported a conclusion that the Monoceros Ring, A13 and TriAnd Ring were stellar overdensities rather kicked out from the main stellar disk, with the velocity dispersion of the RR Lyrae stars found to be higher and consistent with halo membership.
Another 2018 study revealed the very probable presence of disk stars at from the Galactic Center or perhaps even farther, significantly beyond approximately , in which it was once believed to be the abrupt drop-off of the stellar density of the disk, meaning that few or no stars were expected to be above this limit, save for stars that belong to the old population of the galactic halo.
A 2020 study predicted the edge of the Milky Way's dark matter halo being around , which translates to a diameter of . The Milky Way's stellar disk is also estimated to be approximately up to thick.
Mass
The Milky Way is approximately 890 billion to 1.54 trillion times the mass of the Sun in total (8.9 to 1.54 solar masses), although stars and planets make up only a small part of this. Estimates of the mass of the Milky Way vary, depending upon the method and data used. The low end of the estimate range is 5.8 solar masses (), somewhat less than that of the Andromeda Galaxy. Measurements using the Very Long Baseline Array in 2009 found velocities as large as for stars at the outer edge of the Milky Way.
Because the orbital velocity depends on the total mass inside the orbital radius, this suggests that the Milky Way is more massive, roughly equaling the mass of Andromeda Galaxy at 7 within of its center. In 2010, a measurement of the radial velocity of halo stars found that the mass enclosed within 80 kiloparsecs is 7 . In a 2014 study, the mass of the entire Milky Way is estimated to be 8.5 , but this is only half the mass of the Andromeda Galaxy. A recent 2019 mass estimate for the Milky Way is 1.29 .
Much of the mass of the Milky Way seems to be dark matter, an unknown and invisible form of matter that interacts gravitationally with ordinary matter. A dark matter halo is conjectured to spread out relatively uniformly to a distance beyond one hundred kiloparsecs (kpc) from the Galactic Center. Mathematical models of the Milky Way suggest that the mass of dark matter is 1–1.5 . 2013 and 2014 studies indicate a range in mass, as large as 4.5 and as small as 8 . By comparison, the total mass of all the stars in the Milky Way is estimated to be between 4.6 and 6.43 .
In addition to the stars, there is also interstellar gas, comprising 90% hydrogen and 10% helium by mass, with two thirds of the hydrogen found in the atomic form and the remaining one-third as molecular hydrogen. The mass of the Milky Way's interstellar gas is equal to between 10% and 15% of the total mass of its stars. Interstellar dust accounts for an additional 1% of the total mass of the gas.
In March 2019, astronomers reported that the virial mass of the Milky Way Galaxy is 1.54 trillion solar masses within a radius of about , over twice as much as was determined in earlier studies, suggesting that about 90% of the mass of the galaxy is dark matter.
In September 2023, astronomers reported that the virial mass of the Milky Way Galaxy is only 2.06 1011 solar masses, only a 10th of the mass of previous studies. The mass was determined from data of the Gaia spacecraft.
Contents
The Milky Way contains between 100 and 400 billion stars and at least that many planets. An exact figure would depend on counting the number of very-low-mass stars, which are difficult to detect, especially at distances of more than from the Sun. As a comparison, the neighboring Andromeda Galaxy contains an estimated one trillion (1012) stars. The Milky Way may contain ten billion white dwarfs, a billion neutron stars, and a hundred million stellar black holes. Filling the space between the stars is a disk of gas and dust called the interstellar medium. This disk has at least a comparable extent in radius to the stars, whereas the thickness of the gas layer ranges from hundreds of light-years for the colder gas to thousands of light-years for warmer gas.
The disk of stars in the Milky Way does not have a sharp edge beyond which there are no stars. Rather, the concentration of stars decreases with distance from the center of the Milky Way. Beyond a radius of roughly 40,000 light years (13 kpc) from the center, the number of stars per cubic parsec drops much faster with radius. Surrounding the galactic disk is a spherical galactic halo of stars and globular clusters that extends farther outward, but is limited in size by the orbits of two Milky Way satellites, the Large and Small Magellanic Clouds, whose closest approach to the Galactic Center is about . At this distance or beyond, the orbits of most halo objects would be disrupted by the Magellanic Clouds. Hence, such objects would probably be ejected from the vicinity of the Milky Way. The integrated absolute visual magnitude of the Milky Way is estimated to be around −20.9.
Both gravitational microlensing and planetary transit observations indicate that there may be at least as many planets bound to stars as there are stars in the Milky Way, and microlensing measurements indicate that there are more rogue planets not bound to host stars than there are stars. The Milky Way contains an average of at least one planet per star, resulting in 100–400 billion planets, according to a January 2013 study of the five-planet star system Kepler-32 by the Kepler space observatory. A different January 2013 analysis of Kepler data estimated that at least 17 billion Earth-sized exoplanets reside in the Milky Way.
In November 2013, astronomers reported, based on Kepler space mission data, that there could be as many as 40 billion Earth-sized planets orbiting in the habitable zones of Sun-like stars and red dwarfs within the Milky Way. 11 billion of these estimated planets may be orbiting Sun-like stars. The nearest exoplanet may be 4.2 light-years away, orbiting the red dwarf Proxima Centauri, according to a 2016 study. Such Earth-sized planets may be more numerous than gas giants, though harder to detect at great distances given their small size. Besides exoplanets, "exocomets", comets beyond the Solar System, have also been detected and may be common in the Milky Way. More recently, in November 2020, over 300 million habitable exoplanets are estimated to exist in the Milky Way Galaxy.
When compared to other more distant galaxies in the universe, the Milky Way galaxy has a below average amount of neutrino luminosity making our galaxy a "neutrino desert".
Structure
The Milky Way consists of a bar-shaped core region surrounded by a warped disk of gas, dust and stars. The mass distribution within the Milky Way closely resembles the type Sbc in the Hubble classification, which represents spiral galaxies with relatively loosely wound arms. Astronomers first began to conjecture that the Milky Way is a barred spiral galaxy, rather than an ordinary spiral galaxy, in the 1960s. These conjectures were confirmed by the Spitzer Space Telescope observations in 2005 that showed the Milky Way's central bar to be larger than previously thought.
Galactic Center
The Sun is from the Galactic Center. This value is estimated using geometric-based methods or by measuring selected astronomical objects that serve as standard candles, with different techniques yielding various values within this approximate range. In the inner few kiloparsecs (around 10,000 light-years radius) is a dense concentration of mostly old stars in a roughly spheroidal shape called the bulge. It has been proposed that the Milky Way lacks a bulge due to a collision and merger between previous galaxies, and that instead it only has a pseudobulge formed by its central bar. However, confusion in the literature between the (peanut shell)-shaped structure created by instabilities in the bar, versus a possible bulge with an expected half-light radius of 0.5 kpc, abounds.
The Galactic Center is marked by an intense radio source named (pronounced Sagittarius A-star). The motion of material around the center indicates that Sagittarius A* harbors a massive, compact object. This concentration of mass is best explained as a supermassive black hole (SMBH) with an estimated mass of 4.1–4.5 million times the mass of the Sun. The rate of accretion of the SMBH is consistent with an inactive galactic nucleus, being estimated at per year. Observations indicate that there are SMBHs located near the center of most normal galaxies.
The nature of the Milky Way's bar is actively debated, with estimates for its half-length and orientation spanning from and 10–50 degrees relative to the line of sight from Earth to the Galactic Center. Certain authors advocate that the Milky Way features two distinct bars, one nestled within the other. However, RR Lyrae-type stars do not trace a prominent Galactic bar. The bar may be surrounded by a ring called the "5 kpc ring" that contains a large fraction of the molecular hydrogen present in the Milky Way, as well as most of the Milky Way's star formation activity. Viewed from the Andromeda Galaxy, it would be the brightest feature of the Milky Way. X-ray emission from the core is aligned with the massive stars surrounding the central bar and the Galactic ridge.
In June 2023, astronomers led by Naoko Kurahashi Neilson reported using a new cascade neutrino technique to detect, for the first time, the release of neutrinos from the galactic plane of the Milky Way galaxy, creating the first neutrino view of the Milky Way.
Gamma rays and x-rays
Since 1970, various gamma-ray detection missions have discovered 511-keV gamma rays coming from the general direction of the Galactic Center. These gamma rays are produced by positrons (antielectrons) annihilating with electrons. In 2008 it was found that the distribution of the sources of the gamma rays resembles the distribution of low-mass X-ray binaries, seeming to indicate that these X-ray binaries are sending positrons (and electrons) into interstellar space where they slow down and annihilate. The observations were both made by NASA and ESA's satellites. In 1970 gamma ray detectors found that the emitting region was about 10,000 light-years across with a luminosity of about 10,000 suns.
In 2010, two gigantic spherical bubbles of high energy gamma-emission were detected to the north and the south of the Milky Way core, using data from the Fermi Gamma-ray Space Telescope. The diameter of each of the bubbles is about (or about 1/4 of the galaxy's estimated diameter); they stretch up to Grus and to Virgo on the night-sky of the southern hemisphere. Subsequently, observations with the Parkes Telescope at radio frequencies identified polarized emission that is associated with the Fermi bubbles. These observations are best interpreted as a magnetized outflow driven by star formation in the central of the Milky Way.
Later, on January 5, 2015, NASA reported observing an X-ray flare 400 times brighter than usual, a record-breaker, from Sagittarius A*. The unusual event may have been caused by the breaking apart of an asteroid falling into the black hole or by the entanglement of magnetic field lines within gas flowing into Sagittarius A*.
Spiral arms
Outside the gravitational influence of the Galactic bar, the structure of the interstellar medium and stars in the disk of the Milky Way is organized into four spiral arms. Spiral arms typically contain a higher density of interstellar gas and dust than the Galactic average as well as a greater concentration of star formation, as traced by H II regions and molecular clouds.
The Milky Way's spiral structure is uncertain, and there is currently no consensus on the nature of the Milky Way's arms. Perfect logarithmic spiral patterns only crudely describe features near the Sun, because galaxies commonly have arms that branch, merge, twist unexpectedly, and feature a degree of irregularity. The possible scenario of the Sun within a spur / Local arm emphasizes that point and indicates that such features are probably not unique, and exist elsewhere in the Milky Way. Estimates of the pitch angle of the arms range from about 7° to 25°. There are thought to be four spiral arms that all start near the Milky Way Galaxy's center. These are named as follows, with the positions of the arms shown in the image:
Two spiral arms, the Scutum–Centaurus arm and the Carina–Sagittarius arm, have tangent points inside the Sun's orbit about the center of the Milky Way. If these arms contain an overdensity of stars compared to the average density of stars in the Galactic disk, it would be detectable by counting the stars near the tangent point. Two surveys of near-infrared light, which is sensitive primarily to red giants and not affected by dust extinction, detected the predicted overabundance in the Scutum–Centaurus arm but not in the Carina–Sagittarius arm: the Scutum–Centaurus Arm contains approximately 30% more red giants than would be expected in the absence of a spiral arm.
This observation suggests that the Milky Way possesses only two major stellar arms: the Perseus arm and the Scutum–Centaurus arm. The rest of the arms contain excess gas but not excess old stars. In December 2013, astronomers found that the distribution of young stars and star-forming regions matches the four-arm spiral description of the Milky Way. Thus, the Milky Way appears to have two spiral arms as traced by old stars and four spiral arms as traced by gas and young stars. The explanation for this apparent discrepancy is unclear.
The Near 3 kpc Arm (also called the Expanding 3 kpc Arm or simply the 3 kpc Arm) was discovered in the 1950s by astronomer van Woerden and collaborators through 21 centimeter radio measurements of H (atomic hydrogen). It was found to be expanding away from the central bulge at more than 50 km/s. It is located in the fourth galactic quadrant at a distance of about 5.2 kpc from the Sun and 3.3 kpc from the Galactic Center. The Far 3 kpc Arm was discovered in 2008 by astronomer Tom Dame (Center for Astrophysics Harvard & Smithsonian). It is located in the first galactic quadrant at a distance of 3 kpc (about 10,000 ly) from the Galactic Center.
A simulation published in 2011 suggested that the Milky Way may have obtained its spiral arm structure as a result of repeated collisions with the Sagittarius Dwarf Elliptical Galaxy.
It has been suggested that the Milky Way contains two different spiral patterns: an inner one, formed by the Sagittarius arm, that rotates fast and an outer one, formed by the Carina and Perseus arms, whose rotation velocity is slower and whose arms are tightly wound. In this scenario, suggested by numerical simulations of the dynamics of the different spiral arms, the outer pattern would form an outer pseudoring, and the two patterns would be connected by the Cygnus arm.
Outside of the major spiral arms is the Monoceros Ring (or Outer Ring), a ring of gas and stars torn from other galaxies billions of years ago. However, several members of the scientific community recently restated their position affirming the Monoceros structure is nothing more than an over-density produced by the flared and warped thick disk of the Milky Way. The structure of the Milky Way's disk is warped along an "S" curve.
Halo
The Galactic disk is surrounded by a spheroidal halo of old stars and globular clusters, of which 90% lie within of the Galactic Center. However, a few globular clusters have been found farther, such as PAL 4 and AM 1 at more than 200,000 light-years from the Galactic Center. About 40% of the Milky Way's clusters are on retrograde orbits, which means they move in the opposite direction from the Milky Way rotation. The globular clusters can follow rosette orbits about the Milky Way, in contrast to the elliptical orbit of a planet around a star.
Although the disk contains dust that obscures the view in some wavelengths, the halo component does not. Active star formation takes place in the disk (especially in the spiral arms, which represent areas of high density), but does not take place in the halo, as there is little cool gas to collapse into stars. Open clusters are also located primarily in the disk.
Discoveries in the early 21st century have added dimension to the knowledge of the Milky Way's structure. With the discovery that the disk of the Andromeda Galaxy (M31) extends much farther than previously thought, the possibility of the disk of the Milky Way extending farther is apparent, and this is supported by evidence from the discovery of the Outer Arm extension of the Cygnus Arm and of a similar extension of the Scutum–Centaurus Arm. With the discovery of the Sagittarius Dwarf Elliptical Galaxy came the discovery of a ribbon of galactic debris as the polar orbit of the dwarf and its interaction with the Milky Way tears it apart. Similarly, with the discovery of the Canis Major Dwarf Galaxy, it was found that a ring of galactic debris from its interaction with the Milky Way encircles the Galactic disk.
The Sloan Digital Sky Survey of the northern sky shows a huge and diffuse structure (spread out across an area around 5,000 times the size of a full moon) within the Milky Way that does not seem to fit within current models. The collection of stars rises close to perpendicular to the plane of the spiral arms of the Milky Way. The proposed likely interpretation is that a dwarf galaxy is merging with the Milky Way. This galaxy is tentatively named the Virgo Stellar Stream and is found in the direction of Virgo about away.
Gaseous halo
In addition to the stellar halo, the Chandra X-ray Observatory, XMM-Newton, and Suzaku have provided evidence that there is also a gaseous halo containing a large amount of hot gas. This halo extends for hundreds of thousands of light-years, much farther than the stellar halo and close to the distance of the Large and Small Magellanic Clouds. The mass of this hot halo is nearly equivalent to the mass of the Milky Way itself. The temperature of this halo gas is between 1 and 2.5 million K (1.8 and 4.5 million °F).
Observations of distant galaxies indicate that the Universe had about one-sixth as much baryonic (ordinary) matter as dark matter when it was just a few billion years old. However, only about half of those baryons are accounted for in the modern Universe based on observations of nearby galaxies like the Milky Way. If the finding that the mass of the halo is comparable to the mass of the Milky Way is confirmed, it could be the identity of the missing baryons around the Milky Way.
Galactic rotation
The stars and gas in the Milky Way rotate about its center differentially, meaning that the rotation period varies with location. As is typical for spiral galaxies, the orbital speed of most stars in the Milky Way does not depend strongly on their distance from the center. Away from the central bulge or outer rim, the typical stellar orbital speed is between . Hence the orbital period of the typical star is directly proportional only to the length of the path traveled. This is unlike the situation within the Solar System, where two-body gravitational dynamics dominate, and different orbits have significantly different velocities associated with them. The rotation curve (shown in the figure) describes this rotation. Toward the center of the Milky Way the orbit speeds are too low, whereas beyond 7 kpcs the speeds are too high to match what would be expected from the universal law of gravitation.
If the Milky Way contained only the mass observed in stars, gas, and other baryonic (ordinary) matter, the rotational speed would decrease with distance from the center. However, the observed curve is relatively flat, indicating that there is additional mass that cannot be detected directly with electromagnetic radiation. This inconsistency is attributed to dark matter. The rotation curve of the Milky Way agrees with the universal rotation curve of spiral galaxies, the best evidence for the existence of dark matter in galaxies. Alternatively, a minority of astronomers propose that a modification of the law of gravity may explain the observed rotation curve.
Formation
History
The Milky Way began as one or several small overdensities in the mass distribution in the Universe shortly after the Big Bang 13.61 billion years ago. Some of these overdensities were the seeds of globular clusters in which the oldest remaining stars in what is now the Milky Way formed. Nearly half the matter in the Milky Way may have come from other distant galaxies. These stars and clusters now comprise the stellar halo of the Milky Way. Within a few billion years of the birth of the first stars, the mass of the Milky Way was large enough so that it was spinning relatively quickly. Due to conservation of angular momentum, this led the gaseous interstellar medium to collapse from a roughly spheroidal shape to a disk. Therefore, later generations of stars formed in this spiral disk. Most younger stars, including the Sun, are observed to be in the disk.
Since the first stars began to form, the Milky Way has grown through both galaxy mergers (particularly early in the Milky Way's growth) and accretion of gas directly from the Galactic halo. The Milky Way is currently accreting material from several small galaxies, including two of its largest satellite galaxies, the Large and Small Magellanic Clouds, through the Magellanic Stream. Direct accretion of gas is observed in high-velocity clouds like the Smith Cloud.
Cosmological simulations indicate that, 11 billion years ago, it merged with a particularly large galaxy that has been labeled the Kraken. Properties of the Milky Way such as stellar mass, angular momentum, and metallicity in its outermost regions suggest it has undergone no mergers with large galaxies in the last 10 billion years. This lack of recent major mergers is unusual among similar spiral galaxies. Its neighbour the Andromeda Galaxy appears to have a more typical history shaped by more recent mergers with relatively large galaxies.
According to recent studies, the Milky Way as well as the Andromeda Galaxy lie in what in the galaxy color–magnitude diagram is known as the "green valley", a region populated by galaxies in transition from the "blue cloud" (galaxies actively forming new stars) to the "red sequence" (galaxies that lack star formation). Star-formation activity in green valley galaxies is slowing as they run out of star-forming gas in the interstellar medium. In simulated galaxies with similar properties, star formation will typically have been extinguished within about five billion years from now, even accounting for the expected, short-term increase in the rate of star formation due to the collision between both the Milky Way and the Andromeda Galaxy. Measurements of other galaxies similar to the Milky Way suggest it is among the reddest and brightest spiral galaxies that are still forming new stars and it is just slightly bluer than the bluest red sequence galaxies.
Age and cosmological history
Globular clusters are among the oldest objects in the Milky Way, which thus set a lower limit on the age of the Milky Way. The ages of individual stars in the Milky Way can be estimated by measuring the abundance of long-lived radioactive elements such as thorium-232 and uranium-238, then comparing the results to estimates of their original abundance, a technique called nucleocosmochronology. These yield values of about for CS 31082-001 and for BD +17° 3248.
Once a white dwarf is formed, it begins to undergo radiative cooling and the surface temperature steadily drops. By measuring the temperatures of the coolest of these white dwarfs and comparing them to their expected initial temperature, an age estimate can be made. With this technique, the age of the globular cluster M4 was estimated as . Age estimates of the oldest of these clusters gives a best fit estimate of 12.6 billion years, and a 95% confidence upper limit of 16 billion years.
In November 2018, astronomers reported the discovery of one of the oldest stars in the universe. About 13.5 billion-years-old, 2MASS J18082002-5104378 B is a tiny ultra metal-poor (UMP) star made almost entirely of materials released from the Big Bang, and is possibly one of the first stars. The discovery of the star in the Milky Way Galaxy suggests that the galaxy may be at least 3 billion years older than previously thought.
Several individual stars have been found in the Milky Way's halo with measured ages very close to the 13.80-billion-year age of the Universe. In 2007, a star in the galactic halo, HE 1523-0901, was estimated to be about 13.2 billion years old. As the oldest known object in the Milky Way at that time, this measurement placed a lower limit on the age of the Milky Way. This estimate was made using the UV-Visual Echelle Spectrograph of the Very Large Telescope to measure the relative strengths of spectral lines caused by the presence of thorium and other elements created by the R-process. The line strengths yield abundances of different elemental isotopes, from which an estimate of the age of the star can be derived using nucleocosmochronology. Another star, HD 140283, has been estimated at 14.5 ± 0.7 billion years old.
According to observations utilizing adaptive optics to correct for Earth's atmospheric distortion, stars in the galaxy's bulge date to about 12.8 billion years old.
The age of stars in the galactic thin disk has also been estimated using nucleocosmochronology. Measurements of thin disk stars yield an estimate that the thin disk formed 8.8 ± 1.7 billion years ago. These measurements suggest there was a hiatus of almost 5 billion years between the formation of the galactic halo and the thin disk. Recent analysis of the chemical signatures of thousands of stars suggests that stellar formation might have dropped by an order of magnitude at the time of disk formation, 10 to 8 billion years ago, when interstellar gas was too hot to form new stars at the same rate as before.
The satellite galaxies surrounding the Milky Way are not randomly distributed but seem to be the result of a breakup of some larger system producing a ring structure 500,000 light-years in diameter and 50,000 light-years wide. Close encounters between galaxies, like that expected in 4 billion years with the Andromeda Galaxy, can rip off huge tails of gas, which, over time can coalesce to form dwarf galaxies in a ring at an arbitrary angle to the main disc.
Intergalactic neighbourhood
The Milky Way and the Andromeda Galaxy are a binary system of giant spiral galaxies belonging to a group of 50 closely bound galaxies known as the Local Group, surrounded by a Local Void, itself being part of the Local Sheet and in turn the Virgo Supercluster. Surrounding the Virgo Supercluster are a number of voids, devoid of many galaxies, the Microscopium Void to the "north", the Sculptor Void to the "left", the Boötes Void to the "right" and the Canes-Major Void to the "south". These voids change shape over time, creating filamentous structures of galaxies. The Virgo Supercluster, for instance, is being drawn towards the Great Attractor, which in turn forms part of a greater structure, called Laniakea.
Two smaller galaxies and a number of dwarf galaxies in the Local Group orbit the Milky Way. The largest of these is the Large Magellanic Cloud with a diameter of 32,200 light-years. It has a close companion, the Small Magellanic Cloud. The Magellanic Stream is a stream of neutral hydrogen gas extending from these two small galaxies across 100° of the sky. The stream is thought to have been dragged from the Magellanic Clouds in tidal interactions with the Milky Way. Some of the dwarf galaxies orbiting the Milky Way are Canis Major Dwarf (the closest), Sagittarius Dwarf Elliptical Galaxy, Ursa Minor Dwarf, Sculptor Dwarf, Sextans Dwarf, Fornax Dwarf, and Leo I Dwarf.
The smallest dwarf galaxies of the Milky Way are only 500 light-years in diameter. These include Carina Dwarf, Draco Dwarf, and Leo II Dwarf. There may still be undetected dwarf galaxies that are dynamically bound to the Milky Way, which is supported by the detection of nine new satellites of the Milky Way in a relatively small patch of the night sky in 2015. There are some dwarf galaxies that have already been absorbed by the Milky Way, such as the progenitor of Omega Centauri.
In 2005
with further confirmation in 2012 researchers reported that most satellite galaxies of the Milky Way lie in a very large disk and orbit in the same direction. This came as a surprise: according to standard cosmology, the satellite galaxies should form in dark matter halos, and they should be widely distributed and moving in random directions. This discrepancy is still not explained.
In January 2006, researchers reported that the heretofore unexplained warp in the disk of the Milky Way has now been mapped and found to be a ripple or vibration set up by the Large and Small Magellanic Clouds as they orbit the Milky Way, causing vibrations when they pass through its edges. Previously, these two galaxies, at around 2% of the mass of the Milky Way, were considered too small to influence the Milky Way. However, in a computer model, the movement of these two galaxies creates a dark matter wake that amplifies their influence on the larger Milky Way.
Current measurements suggest the Andromeda Galaxy is approaching the Milky Way at . In 4.3 billion years, there may be an Andromeda–Milky Way collision, depending on the importance of unknown lateral components to the galaxies' relative motion. If they collide, the chance of individual stars colliding with each other is extremely low, but instead the two galaxies will merge to form a single elliptical galaxy or perhaps a large disk galaxy over the course of about six billion years.
Velocity
Although special relativity states that there is no "preferred" inertial frame of reference in space with which to compare the Milky Way, the Milky Way does have a velocity with respect to cosmological frames of reference.
One such frame of reference is the Hubble flow, the apparent motions of galaxy clusters due to the expansion of space. Individual galaxies, including the Milky Way, have peculiar velocities relative to the average flow. Thus, to compare the Milky Way to the Hubble flow, one must consider a volume large enough so that the expansion of the Universe dominates over local, random motions. A large enough volume means that the mean motion of galaxies within this volume is equal to the Hubble flow. Astronomers believe the Milky Way is moving at approximately with respect to this local co-moving frame of reference.
The Milky Way is moving in the general direction of the Great Attractor and other galaxy clusters, including the Shapley Supercluster, behind it. The Local Group, a cluster of gravitationally bound galaxies containing, among others, the Milky Way and the Andromeda Galaxy, is part of a supercluster called the Local Supercluster, centered near the Virgo Cluster: although they are moving away from each other at as part of the Hubble flow, this velocity is less than would be expected given the 16.8 million pc distance due to the gravitational attraction between the Local Group and the Virgo Cluster.
Another reference frame is provided by the cosmic microwave background (CMB), in which the CMB temperature is least distorted by Doppler shift (zero dipole moment). The Milky Way is moving at with respect to this frame, toward 10.5 right ascension, −24° declination (J2000 epoch, near the center of Hydra). This motion is observed by satellites such as the Cosmic Background Explorer (COBE) and the Wilkinson Microwave Anisotropy Probe (WMAP) as a dipole contribution to the CMB, as photons in equilibrium in the CMB frame get blue-shifted in the direction of the motion and red-shifted in the opposite direction.
| Physical sciences | Astronomy | null |
2589751 | https://en.wikipedia.org/wiki/Fluorescence%20correlation%20spectroscopy | Fluorescence correlation spectroscopy | Fluorescence correlation spectroscopy (FCS) is a statistical analysis, via time correlation, of stationary fluctuations of the fluorescence intensity. Its theoretical underpinning originated from L. Onsager's regression hypothesis. The analysis provides kinetic parameters of the physical processes underlying the fluctuations. One of the interesting applications of this is an analysis of the concentration fluctuations of fluorescent particles (molecules) in solution. In this application, the fluorescence emitted from a very tiny space in solution containing a small number of fluorescent particles (molecules) is observed. The fluorescence intensity is fluctuating due to Brownian motion of the particles. In other words, the number of the particles in the sub-space defined by the optical system is randomly changing around the average number. The analysis gives the average number of fluorescent particles and average diffusion time, when the particle is passing through the space. Eventually, both the concentration and size of the particle (molecule) are determined. Both parameters are important in biochemical research, biophysics, and chemistry.
FCS is such a sensitive analytical tool because it observes a small number of molecules (nanomolar to picomolar concentrations) in a small volume (~1 μm3). In contrast to other methods (such as HPLC analysis) FCS has no physical separation process; instead, it achieves its spatial resolution through its optics. Furthermore, FCS enables observation of fluorescence-tagged molecules in the biochemical pathway in intact living cells. This opens a new area, "in situ or in vivo biochemistry": tracing the biochemical pathway in intact cells and organs.
Commonly, FCS is employed in the context of optical microscopy, in particular confocal microscopy or two-photon excitation microscopy. In these techniques light is focused on a sample and the measured fluorescence intensity fluctuations (due to diffusion, physical or chemical reactions, aggregation, etc.) are analyzed using the temporal autocorrelation. Because the measured property is essentially related to the magnitude and/or the amount of fluctuations, there is an optimum measurement regime at the level when individual species enter or exit the observation volume (or turn on and off in the volume). When too many entities are measured at the same time the overall fluctuations are small in comparison to the total signal and may not be resolvable – in the other direction, if the individual fluctuation-events are too sparse in time, one measurement may take prohibitively too long. FCS is in a way the fluorescent counterpart to dynamic light scattering, which uses coherent light scattering, instead of (incoherent) fluorescence.
When an appropriate model is known, FCS can be used to obtain quantitative information such as
diffusion coefficients
hydrodynamic radii
average concentrations
kinetic chemical reaction rates
singlet-triplet dynamics
Because fluorescent markers come in a variety of colors and can be specifically bound to a particular molecule (e.g. proteins, polymers, metal-complexes, etc.), it is possible to study the behavior of individual molecules (in rapid succession in composite solutions). With the development of sensitive detectors such as avalanche photodiodes the detection of the fluorescence signal coming from individual molecules in highly dilute samples has become practical. With this emerged the possibility to conduct FCS experiments in a wide variety of specimens, ranging from materials science to biology. The advent of engineered cells with genetically tagged proteins (like green fluorescent protein) has made FCS a common tool for studying molecular dynamics in living cells.
History
Signal-correlation techniques were first experimentally applied to fluorescence in 1972 by Magde, Elson, and Webb, who are therefore commonly credited as the inventors of FCS. The technique was further developed in a group of papers by these and other authors soon after, establishing the theoretical foundations and types of applications.
Around 1990, with the ability of detecting sufficiently small number of fluorescence particles, two issues emerged: A non-Gaussian distribution of the fluorescence intensity and the three-dimensional confocal Measurement Volume of a laser-microscopy system. The former led to an analysis of distributions and moments of the fluorescent signals for extracting molecular information, which eventually became a collection of methods known as Brightness Analyses. See Thompson (1991) for a review of that period.
Beginning in 1993, a number of improvements in the measurement techniques—notably using confocal microscopy, and then two-photon microscopy—to better define the measurement volume and reject background—greatly improved the signal-to-noise ratio and allowed single molecule sensitivity. Since then, there has been a renewed interest in FCS, and as of August 2007 there have been over 3,000 papers using FCS found in Web of Science. See Krichevsky and Bonnet for a review. In addition, there has been a flurry of activity extending FCS in various ways, for instance to laser scanning and spinning-disk confocal microscopy (from a stationary, single point measurement), in using cross-correlation (FCCS) between two fluorescent channels instead of autocorrelation, and in using Förster Resonance Energy Transfer (FRET) instead of fluorescence.
Typical setup
The typical FCS setup consists of a laser line (wavelengths ranging typically from 405–633 nm (cw), and from 690–1100 nm (pulsed)), which is reflected into a microscope objective by a dichroic mirror. The laser beam is focused in the sample, which contains fluorescent particles (molecules) in such high dilution, that only a few are within the focal spot (usually 1–100 molecules in one fL). When the particles cross the focal volume, they fluoresce. This light is collected by the same objective and, because it is red-shifted with respect to the excitation light it passes the dichroic mirror reaching a detector, typically a photomultiplier tube, an avalanche photodiode detector or a superconducting nanowire single-photon detector. The resulting electronic signal can be stored either directly as an intensity versus time trace to be analyzed at a later point, or computed to generate the autocorrelation directly (which requires special acquisition cards). The FCS curve by itself only represents a time-spectrum. Conclusions on physical phenomena have to be extracted from there with appropriate models. The parameters of interest are found after fitting the autocorrelation curve to modeled functional forms.
Measurement volume
The measurement volume is a convolution of illumination (excitation) and detection geometries, which result from the optical elements involved. The resulting volume is described mathematically by the point spread function (or PSF), it is essentially the image of a point source. The PSF is often described as an ellipsoid (with unsharp boundaries) of few hundred nanometers in focus diameter, and almost one micrometer along the optical axis. The shape varies significantly (and has a large impact on the resulting FCS curves) depending on the quality of the optical elements (it is crucial to avoid astigmatism and to check the real shape of the PSF on the instrument). In the case of confocal microscopy, and for small pinholes (around one Airy unit), the PSF is well approximated by Gaussians:
where is the peak intensity, r and z are radial and axial position, and and are the radial and axial radii, and . This Gaussian form is assumed in deriving the functional form of the autocorrelation.
Typically is 200–300 nm, and is 2–6 times larger. One common way of calibrating the measurement volume parameters is to perform FCS on a species with known diffusion coefficient and concentration (see below). Diffusion coefficients for common fluorophores in water are given in a later section.
The Gaussian approximation works to varying degrees depending on the optical details, and corrections can sometimes be applied to offset the errors in approximation.
Autocorrelation function
The (temporal) autocorrelation function is the correlation of a time series with itself shifted by time , as a function of :
where is the deviation from the mean intensity. The normalization (denominator) here is the most commonly used for FCS, because then the correlation at , G(0), is related to the average number of particles in the measurement volume.
As an example, raw FCS data and its autocorrelation for freely diffusing Rhodamine 6G are shown in the figure to the right. The plot on top shows the fluorescent intensity versus time. The intensity fluctuates as Rhodamine 6G moves in and out of the focal volume. In the bottom plot is the autocorrelation on the same data. Information about the diffusion rate and concentration can be obtained using one of the models described below.
For a Gaussian illumination profile , the autocorrelation function is given by the general master formula
where the vector denotes the stochastic displacement in space of a fluorophore after time .
The expression is valid if the average number of fluorophores in the focal volume is low and if dark states, etc., of the fluorophore can be ignored. In particular, no assumption was made on the type of diffusive motion under investigation. The formula allows for an interpretation of as (i) a return probability for small beam parameters and (ii) the moment-generating function of if are varied.
Interpreting the autocorrelation function
To extract quantities of interest, the autocorrelation data can be fitted, typically using a nonlinear least squares algorithm. The fit's functional form depends on the type of dynamics (and the optical geometry in question).
Normal diffusion
The fluorescent particles used in FCS are small and thus experience thermal motions in solution. The simplest FCS experiment is thus normal 3D diffusion, for which the autocorrelation is:
where is the ratio of axial to radial radii of the measurement volume, and is the characteristic residence time. This form was derived assuming a Gaussian measurement volume. Typically, the fit would have three free parameters—G(0), , and —from which the diffusion coefficient and fluorophore concentration can be obtained.
With the normalization used in the previous section, G(0) gives the mean number of diffusers in the volume <N>, or equivalently—with knowledge of the observation volume size—the mean concentration:
where the effective volume is found from integrating the Gaussian form of the measurement volume and is given by:
D gives the diffusion coefficient:
Anomalous diffusion
If the diffusing particles are hindered by obstacles or pushed by a force (molecular motors, flow, etc.) the dynamics is often not sufficiently well-described by the normal diffusion model, where the mean squared displacement (MSD) grows linearly with time. Instead the diffusion may be better described as anomalous diffusion, where the temporal dependence of the MSD is non-linear as in the power-law:
where is an anomalous diffusion coefficient. "Anomalous diffusion" commonly refers only to this very generic model, and not the many other possibilities that might be described as anomalous. Also, a power law is, in a strict sense, the expected form only for a narrow range of rigorously defined systems, for instance when the distribution of obstacles is fractal. Nonetheless a power law can be a useful approximation for a wider range of systems.
The FCS autocorrelation function for anomalous diffusion is:
where the anomalous exponent is the same as above, and becomes a free parameter in the fitting.
Using FCS, the anomalous exponent has been shown to be an indication of the degree of molecular crowding (it is less than one and smaller for greater degrees of crowding).
Polydisperse diffusion
If there are diffusing particles with different sizes (diffusion coefficients), it is common to fit to a function that is the sum of single component forms:
where the sum is over the number different sizes of particle, indexed by i, and gives the weighting, which is related to the quantum yield and concentration of each type. This introduces new parameters, which makes the fitting more difficult as a higher-dimensional space must be searched. Nonlinear least square fitting typically becomes unstable with even a small number of s. A more robust fitting scheme, especially useful for polydisperse samples, is the Maximum Entropy Method.
Diffusion with flow
With diffusion together with a uniform flow with velocity in the lateral direction, the autocorrelation is:
where is the average residence time if there is only a flow (no diffusion).
Chemical relaxation
A wide range of possible FCS experiments involve chemical reactions that continually fluctuate from equilibrium because of thermal motions (and then "relax"). In contrast to diffusion, which is also a relaxation process, the fluctuations cause changes between states of different energies. One very simple system showing chemical relaxation would be a stationary binding site in the measurement volume, where particles only produce signal when bound (e.g. by FRET, or if the diffusion time is much faster than the sampling interval). In this case the autocorrelation is:
where
is the relaxation time and depends on the reaction kinetics (on and off rates), and:
is related to the equilibrium constant K.
Most systems with chemical relaxation also show measurable diffusion as well, and the autocorrelation function will depend on the details of the system. If the diffusion and chemical reaction are decoupled, the combined autocorrelation is the product of the chemical and diffusive autocorrelations.
Triplet state correction
The autocorrelations above assume that the fluctuations are not due to changes in the fluorescent properties of the particles. However, for the majority of (bio)organic fluorophores—e.g. green fluorescent protein, rhodamine, Cy3 and Alexa Fluor dyes—some fraction of illuminated particles are excited to a triplet state (or other non-radiative decaying states) and then do not emit photons for a characteristic relaxation time . Typically is on the order of microseconds, which is usually smaller than the dynamics of interest (e.g. ) but large enough to be measured. A multiplicative term is added to the autocorrelation to account for the triplet state. For normal diffusion:
where is the fraction of particles that have entered the triplet state and is the corresponding triplet state relaxation time. If the dynamics of interest are much slower than the triplet state relaxation, the short time component of the autocorrelation can simply be truncated and the triplet term is unnecessary.
Common fluorescent probes
The fluorescent species used in FCS is typically a biomolecule of interest that has been tagged with a fluorophore (using immunohistochemistry for instance), or is a naked fluorophore that is used to probe some environment of interest (e.g. the cytoskeleton of a cell). The following table gives diffusion coefficients of some common fluorophores in water at room temperature, and their excitation wavelengths.
Variations
FCS almost always refers to the single point, single channel, temporal autocorrelation measurement, although the term "fluorescence correlation spectroscopy" out of its historical scientific context implies no such restriction. FCS has been extended in a number of variations by different researchers, with each extension generating another name (usually an acronym).
Spot variation fluorescence correlation spectroscopy (svFCS)
Whereas FCS is a point measurement providing diffusion time at a given observation volume, svFCS is a technique where the observation spot is varied in order to measure diffusion times at different spot sizes. The relationship between the diffusion time and the spot area is linear and could be plotted in order to decipher the major contribution of confinement. The resulting curve is called the diffusion law.
This technique is used in Biology to study the plasma membrane organization on living cells.
where is the y axis intercept. In case of Brownian diffusion, . In case of a confinement due to isolated domains, whereas in case of isolated domains, .
svFCS studies on living cells and simulation papers
Sampling-Volume-Controlled Fluorescence Correlation Spectroscopy (SVC-FCS):
z-scan FCS
FCS with Nano-apertures: breaking the diffraction barrier
STED-FCS:
Fluorescence cross-correlation spectroscopy (FCCS)
FCS is sometimes used to study molecular interactions using differences in diffusion times (e.g. the product of an association reaction will be larger and thus have larger diffusion times than the reactants individually); however, FCS is relatively insensitive to molecular mass as can be seen from the following equation relating molecular mass to the diffusion time of globular particles (e.g. proteins):
where is the viscosity of the sample and is the molecular mass of the fluorescent species. In practice, the diffusion times need to be sufficiently different—a factor of at least 1.6—which means the molecular masses must differ by a factor of 4. Dual color fluorescence cross-correlation spectroscopy (FCCS) measures interactions by cross-correlating two or more fluorescent channels (one channel for each reactant), which distinguishes interactions more sensitively than FCS, particularly when the mass change in the reaction is small.
Brightness analysis methods
This set of methods include number and brightness (N&B), photon counting histogram (PCH), fluorescence intensity distribution analysis (FIDA), and Cumulant Analysis. and Spatial Intensity Distribution Analysis. Combination of multiple methods is also reported.
Fluorescence cross correlation spectroscopy overcomes the weak dependence of diffusion rate on molecular mass by looking at multicolor coincidence. What about homo-interactions? The solution lies in brightness analysis. These methods use the heterogeneity in the intensity distribution of fluorescence to measure the molecular brightness of different species in a sample. Since dimers will contain twice the number of fluorescent labels as monomers, their molecular brightness will be approximately double that of monomers. As a result, the relative brightness is sensitive a measure of oligomerization. The average molecular brightness () is related to the variance () and the average intensity () as follows:
Here and are the fractional intensity and molecular brightness, respectively, of species .
The brightness analysis method might be employed to study the interactions of biomolecules upon binding a non-fluorescent reactant to a fluorescent one. The complex formation causes a change in brightness intensity due to steric shielding, charge transfer, photoisomerization rate, or a combination of these phenomena enabling distinguishing the reactant from the product.
FRET-FCS
Another FCS based approach to studying molecular interactions uses fluorescence resonance energy transfer (FRET) instead of fluorescence, and is called FRET-FCS. With FRET, there are two types of probes, as with FCCS; however, there is only one channel and light is only detected when the two probes are very close—close enough to ensure an interaction. The FRET signal is weaker than with fluorescence, but has the advantage that there is only signal during a reaction (aside from autofluorescence).
Scanning FCS
In Scanning fluorescence correlation spectroscopy (sFCS) the measurement volume is moved across the sample in a defined way. The introduction of scanning is motivated by its ability to alleviate or remove several distinct problems often encountered in standard FCS, and thus, to extend the range of applicability of fluorescence correlation methods in biological systems.
Some variations of FCS are only applicable to serial scanning laser microscopes. Image Correlation Spectroscopy and its variations all were implemented on a scanning confocal or scanning two photon microscope, but transfer to other microscopes, like a spinning disk confocal microscope. Raster ICS (RICS), and position sensitive FCS (PSFCS) incorporate the time delay between parts of the image scan into the analysis. Also, low-dimensional scans (e.g. a circular ring)—only possible on a scanning system—can access time scales between single point and full image measurements. Scanning path has also been made to adaptively follow particles.
Spinning disk FCS and spatial mapping
Any of the image correlation spectroscopy methods can also be performed on a spinning disk confocal microscope, which in practice can obtain faster imaging speeds compared to a laser scanning confocal microscope. This approach has recently been applied to diffusion in a spatially varying complex environment, producing a pixel resolution map of a diffusion coefficient. The spatial mapping of diffusion with FCS has subsequently been extended to the TIRF system. Spatial mapping of dynamics using correlation techniques had been applied before, but only at sparse points or at coarse resolution.
Image correlation spectroscopy (ICS)
When the motion is slow (in biology, for example, diffusion in a membrane), getting adequate statistics from a single-point FCS experiment may take a prohibitively long time. More data can be obtained by performing the experiment in multiple spatial points in parallel, using a laser scanning confocal microscope. This approach has been called Image Correlation Spectroscopy (ICS). The measurements can then be averaged together.
Another variation of ICS performs a spatial autocorrelation on images, which gives information about the concentration of particles. The correlation is then averaged in time. While camera white noise does not autocorrelate over time, it does over space - this creates a white noise amplitude in the spatial autocorrelation function which must be accounted for when fitting the autocorrelation amplitude in order to find the concentration of fluorescent molecules.
A natural extension of the temporal and spatial correlation versions is spatio-temporal ICS (STICS). In STICS there is no explicit averaging in space or time (only the averaging inherent in correlation). In systems with non-isotropic motion (e.g. directed flow, asymmetric diffusion), STICS can extract the directional information. A variation that is closely related to STICS (by the Fourier transform) is k-space Image Correlation Spectroscopy (kICS).
There are cross-correlation versions of ICS as well, which can yield the concentration, distribution and dynamics of co-localized fluorescent molecules. Molecules are considered co-localized when individual fluorescence contributions are indistinguishable due to overlapping point-spread functions of fluorescence intensities.
Particle image correlation spectroscopy (PICS)
Source:
PICS is a powerful analysis tool that resolves correlations on the nanometer length and millisecond timescale. Adapted from methods of spatio-temporal image correlation spectroscopy, it exploits the high positional accuracy of single-particle tracking. While conventional tracking methods break down if multiple particle trajectories intersect, this method works in principle for arbitrarily large molecule densities and dynamical parameters (e.g. diffusion coefficients, velocities) as long as individual molecules can be identified. It is computationally cheap and robust and allows one to identify and quantify motions (e.g. diffusion, active transport, confined diffusion) within an ensemble of particles, without any a priori knowledge about the dynamics.
A particle image cross-correlation spectroscopy (PICCS) extension is available for biological processes that involve multiple interaction partners, as can observed by two-color microscopy.
FCS Super-resolution Optical Fluctuation Imaging (fcsSOFI)
Super-resolution optical fluctuation imaging (SOFI) is a super-resolution technique that achieves spatial resolutions below the diffraction limit by post-processing analysis with correlation equations, similar to FCS. While original reports of SOFI used fluctuations from stationary, blinking of fluorophores, FCS has been combined with SOFI where fluctuations are produced from diffusing probes to produce super-resolution spatial maps of diffusion coefficients. This has been applied to understand diffusion and spatial properties of porous and confined materials. This includes agarose and temperature-responsive PNIPAM hydrogels, liquid crystals, and phase-separated polymers and RNA/protein condensates.
Total internal reflection FCS
Total internal reflection fluorescence (TIRF) is a microscopy approach that is only sensitive to a thin layer near the surface of a coverslip, which greatly minimizes background fluorescence. FCS has been extended to that type of microscope, and is called TIR-FCS. Because the fluorescence intensity in TIRF falls off exponentially with distance from the coverslip (instead of as a Gaussian with a confocal), the autocorrelation function is different.
FCS imaging using Light sheet fluorescence microscopy
Light sheet fluorescence microscopy or selective plane imaging microscopy (SPIM) uses illumination that is done perpendicularly to the direction of observation, by using a thin sheet of (laser) light. Under certain conditions, this illumination principle can be combined with fluorescence correlation spectroscopy, to allow spatially resolved imaging of the mobility and interactions of fluorescing particles such as GFP labelled proteins inside living biological samples.
Other fluorescent dynamical approaches
There are two main non-correlation alternatives to FCS that are widely used to study the dynamics of fluorescent species.
Fluorescence recovery after photobleaching (FRAP)
In FRAP, a region is briefly exposed to intense light, irrecoverably photobleaching fluorophores, and the fluorescence recovery due to diffusion of nearby (non-bleached) fluorophores is imaged. A primary advantage of FRAP over FCS is the ease of interpreting qualitative experiments common in cell biology. Differences between cell lines, or regions of a cell, or before and after application of drug, can often be characterized by simple inspection of movies. FCS experiments require a level of processing and are more sensitive to potentially confounding influences like: rotational diffusion, vibrations, photobleaching, dependence on illumination and fluorescence color, inadequate statistics, etc. It is much easier to change the measurement volume in FRAP, which allows greater control. In practice, the volumes are typically larger than in FCS. While FRAP experiments are typically more qualitative, some researchers are studying FRAP quantitatively and including binding dynamics. A disadvantage of FRAP in cell biology is the free radical perturbation of the cell caused by the photobleaching. It is also less versatile, as it cannot measure concentration or rotational diffusion, or co-localization. FRAP requires a significantly higher concentration of fluorophores than FCS.
Particle tracking
In particle tracking, the trajectories of a set of particles are measured, typically by applying particle tracking algorithms to movies. Particle tracking has the advantage that all the dynamical information is maintained in the measurement, unlike FCS where correlation averages the dynamics to a single smooth curve. The advantage is apparent in systems showing complex diffusion, where directly computing the mean squared displacement allows straightforward comparison to normal or power law diffusion. To apply particle tracking, the particles have to be distinguishable and thus at lower concentration than required of FCS. Also, particle tracking is more sensitive to noise, which can sometimes affect the results unpredictably.
Auto-fluorescence correlation spectroscopy
Recent advances in ultraviolet nanophotonics has led to development of single molecule study on label-free protein by exciting them with deep ultraviolet light and studying the dynamic processes.
Two- and three-photon FCS excitation
Several advantages in both spatial resolution and minimizing photodamage/photobleaching in organic and/or biological samples are obtained by two-photon or three-photon excitation FCS.
| Physical sciences | Spectroscopy | Chemistry |
2590442 | https://en.wikipedia.org/wiki/Square%20degree | Square degree |
A square degree (deg2) is a non-SI unit measure of solid angle. Other denotations include sq. deg. and (°)2. Just as degrees are used to measure parts of a circle, square degrees are used to measure parts of a sphere. Analogous to one degree being equal to radians, a square degree is equal to ()2 steradians (sr), or about sr or about .
The whole sphere has a solid angle of which is approximately :
Examples
The full moon covers only about of the sky when viewed from the surface of the Earth. The Moon is only a half degree across (i.e. a circular diameter of roughly ), so the moon's disk covers a circular area of: ()2, or 0.2 square degrees. The moon varies from 0.188 to depending on its distance from the Earth.
Viewed from Earth, the Sun is roughly half a degree across (the same as the full moon) and covers only as well.
It would take times the full moon (or the Sun) to cover the entire celestial sphere.
Conversely, an average full moon (or the Sun) covers a 2 / fraction, or less than 1/1000 of a percent () of the celestial hemisphere, or above-the-horizon sky.
Assuming the Earth to be a sphere with a surface area of 510 million km2, the area of Northern Ireland () represents a solid angle of , Connecticut () represents a solid angle of , Equatorial Guinea () represents a solid angle of .
The largest constellation, Hydra, covers a solid angle of , whereas the smallest, Crux, covers only .
| Physical sciences | Solid angle | Basics and measurement |
2592906 | https://en.wikipedia.org/wiki/Planetary%20habitability | Planetary habitability | Planetary habitability is the measure of a planet's or a natural satellite's potential to develop and maintain an environment hospitable to life. Life may be generated directly on a planet or satellite endogenously. Research suggests that life may also be transferred from one body to another, through a hypothetical process known as panspermia. Environments do not need to contain life to be considered habitable nor are accepted habitable zones (HZ) the only areas in which life might arise.
As the existence of life beyond Earth is unknown, planetary habitability is largely an extrapolation of conditions on Earth and the characteristics of the Sun and Solar System which appear favorable to life's flourishing. Of particular interest are those factors that have sustained complex, multicellular organisms on Earth and not just simpler, unicellular creatures. Research and theory in this regard is a component of a number of natural sciences, such as astronomy, planetary science and the emerging discipline of astrobiology.
An absolute requirement for life is an energy source, and the notion of planetary habitability implies that many other geophysical, geochemical, and astrophysical criteria must be met before an astronomical body can support life. In its astrobiology roadmap, NASA has defined the principal habitability criteria as "extended regions of liquid water, conditions favorable for the assembly of complex organic molecules, and energy sources to sustain metabolism". In August 2018, researchers reported that water worlds could support life.
Habitability indicators and biosignatures must be interpreted within a planetary and environmental context. In determining the habitability potential of a body, studies focus on its bulk composition, orbital properties, atmosphere, and potential chemical interactions. Stellar characteristics of importance include mass and luminosity, stable variability, and high metallicity. Rocky, wet terrestrial-type planets and moons with the potential for Earth-like chemistry are a primary focus of astrobiological research, although more speculative habitability theories occasionally examine alternative biochemistries and other types of astronomical bodies.
Background
The idea that planets beyond Earth might host life is an ancient one, though historically it was framed by philosophy as much as physical science. The late 20th century saw two breakthroughs in the field. The observation and robotic spacecraft exploration of other planets and moons within the Solar System has provided critical information on defining habitability criteria and allowed for substantial geophysical comparisons between the Earth and other bodies. The discovery of exoplanets, beginning in the early 1990s and accelerating thereafter, has provided further information for the study of possible extraterrestrial life. These findings confirm that the Sun is not unique among stars in hosting planets and expands the habitability research horizon beyond the Solar System.
While Earth is the only place in the Universe known to harbor life, estimates of habitable zones around other stars, along with the discovery of thousands of exoplanets and new insights into the extreme habitats on Earth where organisms known as extremophiles live, suggest that there may be many more habitable places in the Universe than considered possible until very recently. On 4 November 2013, astronomers reported, based on Kepler space mission data, that there could be as many as 40 billion Earth-sized planets orbiting in the habitable zones of Sun-like stars and red dwarfs within the Milky Way. 11 billion of these estimated planets may be orbiting Sun-like stars. The nearest such planet may be 12 light-years away, according to the scientists. As of June 2021, a total of 59 potentially habitable exoplanets have been found.
Stellar characteristics
An understanding of planetary habitability begins with the host star. The classical habitable zone (HZ) is defined for surface conditions only; but a metabolism that does not depend on the stellar light can still exist outside the HZ, thriving in the interior of the planet where liquid water is available.
Under the auspices of SETI's Project Phoenix, scientists Margaret Turnbull and Jill Tarter developed the "HabCat" (or Catalogue of Habitable Stellar Systems) in 2002. The catalogue was formed by winnowing the nearly 120,000 stars of the larger Hipparcos Catalogue into a core group of 17,000 potentially habitable stars, and the selection criteria that were used provide a good starting point for understanding which astrophysical factors are necessary for habitable planets. According to research published in August 2015, very large galaxies may be more favorable to the formation and development of habitable planets than smaller galaxies, like the Milky Way galaxy.
However, what makes a planet habitable is a much more complex question than having a planet located at the right distance from its host star so that water can be liquid on its surface: various geophysical and geodynamical aspects, the radiation, and the host star's plasma environment can influence the evolution of planets and life, if it originated. Liquid water is a necessary but not sufficient condition for life as we know it, as habitability is a function of a multitude of environmental parameters.
Spectral class
The spectral class of a star indicates its photospheric temperature, which (for main-sequence stars) correlates to overall mass. The appropriate spectral range for habitable stars is considered to be "late F" or "G", to "mid-K". This corresponds to temperatures of a little more than 7,000 K down to a little less than 4,000 K (6,700 °C to 3,700 °C); the Sun, a G2 star at 5,777 K, is well within these bounds. This spectral range probably accounts for between 5% and 10% of stars in the local Milky Way galaxy. "Middle-class" stars of this sort have a number of characteristics considered important to planetary habitability:
They live at least a few hundred million years, allowing life a chance to evolve. More luminous main-sequence stars of the "O" classes and many members of the "B" classes usually live less than 500 million years and in exceptional cases less than 10 million.
They emit enough high-frequency ultraviolet radiation to trigger important atmospheric dynamics such as ozone formation, but not so much that ionisation destroys incipient life.
They emit sufficient radiation at wavelengths conducive to photosynthesis.
Liquid water may exist on the surface of planets orbiting them at a distance that does not induce tidal locking.
K-type stars may be able to support life far longer than the Sun.
Whether fainter late K and M class red dwarf stars are also suitable hosts for habitable planets is perhaps the most important open question in the entire field of planetary habitability given their prevalence (habitability of red dwarf systems). Gliese 581 c, a "super-Earth", has been found orbiting in the "habitable zone" (HZ) of a red dwarf and may possess liquid water. However it is also possible that a greenhouse effect may render it too hot to support life, while its neighbor, Gliese 581 d, may be a more likely candidate for habitability. In September 2010, the discovery was announced of another planet, Gliese 581 g, in an orbit between these two planets. However, reviews of the discovery have placed the existence of this planet in doubt, and it is listed as "unconfirmed". In September 2012, the discovery of two planets orbiting Gliese 163 was announced. One of the planets, Gliese 163 c, about 6.9 times the mass of Earth and somewhat hotter, was considered to be within the habitable zone.
A recent study suggests that cooler stars that emit more light in the infrared and near infrared may actually host warmer planets with less ice and incidence of snowball states. These wavelengths are absorbed by their planets' ice and greenhouse gases and remain warmer.
A 2020 study found that about half of Sun-like stars could host rocky, potentially habitable planets. Specifically, they estimated with that, on average, the nearest habitable zone planet around G and K-type stars is about 6 parsecs away, and there are about 4 rocky planets around G and K-type stars within 10 parsecs (32.6 light years) of the Sun.
Stable habitable zone
The habitable zone (HZ) is a shell-shaped region of space surrounding a star in which a planet could maintain liquid water on its surface. The concept was first proposed by astrophysicist Su-Shu Huang in 1959, based on climatic constraints imposed by the host star. After an energy source, liquid water is widely considered the most important ingredient for life, considering how integral it is to all life systems on Earth. However, if life is discovered in the absence of water, the definition of an HZ may have to be greatly expanded.
The inner edge of the HZ is the distance where runaway greenhouse effect vaporize the whole water reservoir and, as a second effect, induce the photodissociation of water vapor and the loss of hydrogen to space. The outer edge of the HZ is the distance from the star where a maximum greenhouse effect fails to keep the surface of the planet above the freezing point, and by (carbon dioxide) condensation.
A "stable" HZ implies two factors. First, the range of an HZ should not vary greatly over time. All stars increase in luminosity as they age, and a given HZ thus migrates outwards, but if this happens too quickly (for example, with a super-massive star) planets may only have a brief window inside the HZ and a correspondingly smaller chance of developing life. Calculating an HZ range and its long-term movement is never straightforward, as negative feedback loops such as the CNO cycle will tend to offset the increases in luminosity. Assumptions made about atmospheric conditions and geology thus have as great an impact on a putative HZ range as does stellar evolution: the proposed parameters of the Sun's HZ, for example, have fluctuated greatly.
Second, no large-mass body such as a gas giant should be present in or relatively close to the HZ, thus disrupting the formation of Earth-size bodies. The matter in the asteroid belt, for example, appears to have been unable to accrete into a planet due to orbital resonances with Jupiter; if the giant had appeared in the region that is now between the orbits of Venus and Mars, Earth would almost certainly not have developed in its present form. However a gas giant inside the HZ might have habitable moons under the right conditions.
Low stellar variation
Changes in luminosity are common to all stars, but the severity of such fluctuations covers a broad range. Most stars are relatively stable, but a significant minority of variable stars often undergo sudden and intense increases in luminosity and consequently in the amount of energy radiated toward bodies in orbit. These stars are considered poor candidates for hosting life-bearing planets, as their unpredictability and energy output changes would negatively impact organisms: living things adapted to a specific temperature range could not survive too great a temperature variation. Further, upswings in luminosity are generally accompanied by massive doses of gamma ray and X-ray radiation which might prove lethal. Atmospheres do mitigate such effects, but their atmosphere might not be retained by planets orbiting variables, because the high-frequency energy buffeting these planets would continually strip them of their protective covering.
The Sun, in this respect as in many others, is relatively benign: the variation between its maximum and minimum energy output is roughly 0.1% over its 11-year solar cycle. There is strong (though not undisputed) evidence that even minor changes in the Sun's luminosity have had significant effects on the Earth's climate well within the historical era: the Little Ice Age of the mid-second millennium, for instance, may have been caused by a relatively long-term decline in the Sun's luminosity. Thus, a star does not have to be a true variable for differences in luminosity to affect habitability. Of known solar analogs, one that closely resembles the Sun is considered to be 18 Scorpii; unfortunately for the prospects of life existing in its proximity, the only significant difference between the two bodies is the amplitude of the solar cycle, which appears to be much greater for 18 Scorpii.
High metallicity
While the bulk of material in any star is hydrogen and helium, there is a significant variation in the amount of heavier elements (metals). A high proportion of metals in a star correlates to the amount of heavy material initially available in the protoplanetary disk. A smaller amount of metal makes the formation of planets much less likely, under the solar nebula theory of planetary system formation. Any planets that did form around a metal-poor star would probably be low in mass, and thus unfavorable for life. Spectroscopic studies of systems where exoplanets have been found to date confirm the relationship between high metal content and planet formation: "Stars with planets, or at least with planets similar to the ones we are finding today, are clearly more metal rich than stars without planetary companions." This relationship between high metallicity and planet formation also means that habitable systems are more likely to be found around stars of younger generations, since stars that formed early in the universe's history have low metal content.
Planetary characteristics
Habitability indicators and biosignatures must be interpreted within a planetary and environmental context. Whether a planet will emerge as habitable depends on the sequence of events that led to its formation, which could include the production of organic molecules in molecular clouds and protoplanetary disks, delivery of materials during and after planetary accretion, and the orbital location in the planetary system. The chief assumption about habitable planets is that they are terrestrial. Such planets, roughly within one order of magnitude of Earth mass, are primarily composed of silicate rocks, and have not accreted the gaseous outer layers of hydrogen and helium found on gas giants. The possibility that life could evolve in the cloud tops of giant planets has not been decisively ruled out, though it is considered unlikely, as they have no surface and their gravity is enormous. The natural satellites of giant planets, meanwhile, remain valid candidates for hosting life.
In February 2011 the Kepler Space Observatory Mission team released a list of 1235 extrasolar planet candidates, including 54 that may be in the habitable zone. Six of the candidates in this zone are smaller than twice the size of Earth. A more recent study found that one of these candidates (KOI 326.01) is much larger and hotter than first reported. Based on the findings, the Kepler team estimated there to be "at least 50 billion planets in the Milky Way" of which "at least 500 million" are in the habitable zone.
In analyzing which environments are likely to support life, a distinction is usually made between simple, unicellular organisms such as bacteria and archaea and complex metazoans (animals). Unicellularity necessarily precedes multicellularity in any hypothetical tree of life, and where single-celled organisms do emerge there is no assurance that greater complexity will then develop. The planetary characteristics listed below are considered crucial for life generally, but in every case multicellular organisms are more picky than unicellular life.
In August 2021, a new class of habitable planets, named ocean planets, which involves "hot, ocean-covered planets with hydrogen-rich atmospheres", has been reported. Hycean planets may soon be studied for biosignatures by terrestrial telescopes as well as space telescopes, such as the James Webb Space Telescope (JWST), which was launched on 25 December 2021.
Mass and size
Low-mass planets are poor candidates for life for two reasons. First, their lesser gravity makes atmosphere retention difficult. Constituent molecules are more likely to reach escape velocity and be lost to space when buffeted by solar wind or stirred by collision. Planets without a thick atmosphere lack the matter necessary for primal biochemistry, have little insulation and poor heat transfer across their surfaces (for example, Mars, with its thin atmosphere, is colder than the Earth would be if it were at a similar distance from the Sun), and provide less protection against meteoroids and high-frequency radiation. Further, where an atmosphere is less dense than 0.006 Earth atmospheres, water cannot exist in liquid form as the required atmospheric pressure, 4.56 mm Hg (608 Pa) (0.18 inch Hg), does not occur. In addition, a lessened pressure reduces the range of temperatures at which water is liquid.
Secondly, smaller planets have smaller diameters and thus higher surface-to-volume ratios than their larger cousins. Such bodies tend to lose the energy left over from their formation quickly and end up geologically dead, lacking the volcanoes, earthquakes and tectonic activity which supply the surface with life-sustaining material and the atmosphere with temperature moderators like carbon dioxide. Plate tectonics appear particularly crucial, at least on Earth: not only does the process recycle important chemicals and minerals, it also fosters bio-diversity through continent creation and increased environmental complexity and helps create the convective cells necessary to generate Earth's magnetic field. Although geologically active planets with volcanism but no plate tectonics, called Ignan Earths, could also be habitable.
"Low mass" is partly a relative label: the Earth is low mass when compared to the Solar System's gas giants, but it is the largest, by diameter and mass, and the densest of all terrestrial bodies. It is large enough to retain an atmosphere through gravity alone and large enough that its molten core remains a heat engine, driving the diverse geology of the surface (the decay of radioactive elements within a planet's core is the other significant component of planetary heating). Mars, by contrast, is nearly (or perhaps totally) geologically dead and has lost much of its atmosphere. Thus it would be fair to infer that the lower mass limit for habitability lies somewhere between that of Mars and that of Earth or Venus: 0.3 Earth masses has been offered as a rough dividing line for habitable planets. However, a 2008 study by the Harvard-Smithsonian Center for Astrophysics suggests that the dividing line may be higher. Earth may in fact lie on the lower boundary of habitability: if it were any smaller, plate tectonics would be impossible. Venus, which has 85% of Earth's mass, shows no signs of tectonic activity. Conversely, "super-Earths", terrestrial planets with higher masses than Earth, would have higher levels of plate tectonics and thus be firmly placed in the habitable range.
Exceptional circumstances do offer exceptional cases: Jupiter's moon Io (which is smaller than any of the terrestrial planets) is volcanically dynamic because of the gravitational stresses induced by its orbit, and its neighbor Europa may have a liquid ocean or icy slush underneath a frozen shell also due to power generated from orbiting a gas giant.
Saturn's Titan, meanwhile, has an outside chance of harbouring life, as it has retained a thick atmosphere and has liquid methane seas on its surface. Organic-chemical reactions that only require minimum energy are possible in these seas, but whether any living system can be based on such minimal reactions is unclear, and would seem unlikely. These satellites are exceptions, but they prove that mass, as a criterion for habitability, cannot necessarily be considered definitive at this stage of our understanding.
A larger planet is likely to have a more massive atmosphere. A combination of higher escape velocity to retain lighter atoms, and extensive outgassing from enhanced plate tectonics may greatly increase the atmospheric pressure and temperature at the surface compared to Earth. The enhanced greenhouse effect of such a heavy atmosphere would tend to suggest that the habitable zone should be further out from the central star for such massive planets.
Finally, a larger planet is likely to have a large iron core. This allows for a magnetic field to protect the planet from stellar wind and cosmic radiation, which otherwise would tend to strip away planetary atmosphere and to bombard living things with ionized particles. Mass is not the only criterion for producing a magnetic field—as the planet must also rotate fast enough to produce a dynamo effect within its core—but it is a significant component of the process.
The mass of a potentially habitable exoplanet is between 0.1 and 5.0 Earth masses. However it is possible for a habitable world to have a mass as low as 0.0268 Earth Masses. The radius of a potentially habitable exoplanet would range between 0.5 and 1.5 Earth radii.
Orbit and rotation
As with other criteria, stability is the critical consideration in evaluating the effect of orbital and rotational characteristics on planetary habitability. Orbital eccentricity is the difference between a planet's farthest and closest approach to its parent star divided by the sum of said distances. It is a ratio describing the shape of the elliptical orbit. The greater the eccentricity the greater the temperature fluctuation on a planet's surface. Although they are adaptive, living organisms can stand only so much variation, particularly if the fluctuations overlap both the freezing point and boiling point of the planet's main biotic solvent (e.g., water on Earth). If, for example, Earth's oceans were alternately boiling and freezing solid, it is difficult to imagine life as we know it having evolved. The more complex the organism, the greater the temperature sensitivity. The Earth's orbit is almost perfectly circular, with an eccentricity of less than 0.02; other planets in the Solar System (with the exception of Mercury) have eccentricities that are similarly benign.
Habitability is also influenced by the architecture of the planetary system around a star. The evolution and stability of these systems are determined by gravitational dynamics, which drive the orbital evolution of terrestrial planets. Data collected on the orbital eccentricities of extrasolar planets has surprised most researchers: 90% have an orbital eccentricity greater than that found within the Solar System, and the average is fully 0.25. This means that the vast majority of planets have highly eccentric orbits and of these, even if their average distance from their star is deemed to be within the HZ, they nonetheless would be spending only a small portion of their time within the zone.
A planet's movement around its rotational axis must also meet certain criteria if life is to have the opportunity to evolve. A first assumption is that the planet should have moderate seasons. If there is little or no axial tilt (or obliquity) relative to the perpendicular of the ecliptic, seasons will not occur and a main stimulant to biospheric dynamism will disappear. The planet would also be colder than it would be with a significant tilt: when the greatest intensity of radiation is always within a few degrees of the equator, warm weather cannot move poleward and a planet's climate becomes dominated by colder polar weather systems.
If a planet is radically tilted, seasons will be extreme and make it more difficult for a biosphere to achieve homeostasis. The axial tilt of the Earth is higher now (in the Quaternary) than it has been in the past, coinciding with reduced polar ice, warmer temperatures and less seasonal variation. Scientists do not know whether this trend will continue indefinitely with further increases in axial tilt (see Snowball Earth).
The exact effects of these changes can only be computer modelled at present, and studies have shown that even extreme tilts of up to 85 degrees do not absolutely preclude life "provided it does not occupy continental surfaces plagued seasonally by the highest temperature." Not only the mean axial tilt, but also its variation over time must be considered. The Earth's tilt varies between 21.5 and 24.5 degrees over 41,000 years. A more drastic variation, or a much shorter periodicity, would induce climatic effects such as variations in seasonal severity.
Other orbital considerations include:
The planet should rotate relatively quickly so that the day-night cycle is not overlong. If a day takes years, the temperature differential between the day and night side will be pronounced, and problems similar to those noted with extreme orbital eccentricity will come to the fore.
The planet also should rotate quickly enough so that a magnetic dynamo may be started in its iron core to produce a magnetic field.
Change in the direction of the axis rotation (precession) should not be pronounced. In itself, precession need not affect habitability as it changes the direction of the tilt, not its degree. However, precession tends to accentuate variations caused by other orbital deviations; see Milankovitch cycles. Precession on Earth occurs over a 26,000-year cycle.
The Earth's Moon appears to play a crucial role in moderating the Earth's climate by stabilising the axial tilt. It has been suggested that a chaotic tilt may be a "deal-breaker" in terms of habitability—i.e. a satellite the size of the Moon is not only helpful but required to produce stability. This position remains controversial.
In the case of the Earth, the sole Moon is sufficiently massive and orbits so as to significantly contribute to ocean tides, which in turn aids the dynamic churning of Earth's large liquid water oceans. These lunar forces not only help ensure that the oceans do not stagnate, but also play a critical role in Earth's dynamic climate.
Geology
Concentrations of radionuclides in rocky planet mantles may be critical for the habitability of Earth-like planets. Such planets with higher abundances likely lack a persistent dynamo for a significant fraction of their lifetimes, and those with lower concentrations may often be geologically inert. Planetary dynamos create strong magnetic fields which may often be necessary for life to develop or persist as they shield planets from solar winds and cosmic radiation. The electromagnetic emission spectra of stars could be used to identify those which are more likely to host habitable Earth-like planets. As of 2020, radionuclides are thought to be produced by rare stellar processes such as neutron star mergers.
Additional geological characteristics may be essential or major factors in the habitability of natural celestial bodies – including some that may shape the body's heat and magnetic field. Some of these are unknown or not well understood and being investigated by planetary scientists, geochemists and others.
Geochemistry
It is generally assumed that any extraterrestrial life that might exist will be based on the same fundamental biochemistry as found on Earth, as the four elements most vital for life, carbon, hydrogen, oxygen, and nitrogen, are also the most common chemically reactive elements in the universe. Indeed, simple biogenic compounds, such as very simple amino acids such as glycine, have been found in meteorites and in the interstellar medium. These four elements together comprise over 96% of Earth's collective biomass. Carbon has an unparalleled ability to bond with itself and to form a massive array of intricate and varied structures, making it an ideal material for the complex mechanisms that form living cells. Hydrogen and oxygen, in the form of water, compose the solvent in which biological processes take place and in which the first reactions occurred that led to life's emergence. The energy released in the formation of powerful covalent bonds between carbon and oxygen, available by oxidizing organic compounds, is the fuel of all complex life-forms. These four elements together make up amino acids, which in turn are the building blocks of proteins, the substance of living tissue. In addition, neither sulfur (required for the building of proteins) nor phosphorus (needed for the formation of DNA, RNA, and the adenosine phosphates essential to metabolism) are rare.
Relative abundance in space does not always mirror differentiated abundance within planets; of the four life elements, for instance, only oxygen is present in any abundance in the Earth's crust. This can be partly explained by the fact that many of these elements, such as hydrogen and nitrogen, along with their simplest and most common compounds, such as carbon dioxide, carbon monoxide, methane, ammonia, and water, are gaseous at warm temperatures. In the hot region close to the Sun, these volatile compounds could not have played a significant role in the planets' geological formation. Instead, they were trapped as gases underneath the newly formed crusts, which were largely made of rocky, involatile compounds such as silica (a compound of silicon and oxygen, accounting for oxygen's relative abundance). Outgassing of volatile compounds through the first volcanoes would have contributed to the formation of the planets' atmospheres. The Miller–Urey experiment showed that, with the application of energy, simple inorganic compounds exposed to a primordial atmosphere can react to synthesize amino acids.
Even so, volcanic outgassing could not have accounted for the amount of water in Earth's oceans. The vast majority of the water—and arguably carbon—necessary for life must have come from the outer Solar System, away from the Sun's heat, where it could remain solid. Comets impacting with the Earth in the Solar System's early years would have deposited vast amounts of water, along with the other volatile compounds life requires, onto the early Earth, providing a kick-start to the origin of life.
Thus, while there is reason to suspect that the four "life elements" ought to be readily available elsewhere, a habitable system probably also requires a supply of long-term orbiting bodies to seed inner planets. Without comets there is a possibility that life as we know it would not exist on Earth.
Microenvironments and extremophiles
One important qualification to habitability criteria is that only a tiny portion of a planet is required to support life, a so-called Goldilocks Edge or Great Prebiotic Spot. Astrobiologists often concern themselves with "micro-environments", noting that "we lack a fundamental understanding of how evolutionary forces, such as mutation, selection, and genetic drift, operate in micro-organisms that act on and respond to changing micro-environments." Extremophiles are Earth organisms that live in niche environments under severe conditions generally considered inimical to life. Usually (although not always) unicellular, extremophiles include acutely alkaliphilic and acidophilic organisms and others that can survive water temperatures above 100 °C in hydrothermal vents.
The discovery of life in extreme conditions has complicated definitions of habitability, but also generated much excitement amongst researchers in greatly broadening the known range of conditions under which life can persist. For example, a planet that might otherwise be unable to support an atmosphere given the solar conditions in its vicinity, might be able to do so within a deep shadowed rift or volcanic cave. Similarly, craterous terrain might offer a refuge for primitive life. The Lawn Hill crater has been studied as an astrobiological analog, with researchers suggesting rapid sediment infill created a protected microenvironment for microbial organisms; similar conditions may have occurred over the geological history of Mars.
Earth environments that cannot support life are still instructive to astrobiologists in defining the limits of what organisms can endure. The heart of the Atacama Desert, generally considered the driest place on Earth, appears unable to support life, and it has been subject to study by NASA and ESA for that reason: it provides a Mars analog and the moisture gradients along its edges are ideal for studying the boundary between sterility and habitability. The Atacama was the subject of study in 2003 that partly replicated experiments from the Viking landings on Mars in the 1970s; no DNA could be recovered from two soil samples, and incubation experiments were also negative for biosignatures.
Ecological factors
The two current ecological approaches for predicting the potential habitability use 19 or 20 environmental factors, with emphasis on water availability, temperature, presence of nutrients, an energy source, and protection from solar ultraviolet and galactic cosmic radiation.
Classification terminology
The Habitable Exoplanets Catalog uses estimated surface temperature range to classify exoplanets:
hypopsychroplanets – very cold (<−50 °C)
psychroplanets – cold (<−50 to 0 °C)
mesoplanets – medium temperature (0–50 °C; not to be confused with the other definition of mesoplanets)
thermoplanets – hot (50–100 °C)
hyperthermoplanets – (> 100 °C)
Mesoplanets would be ideal for complex life, whereas hypopsychroplanets and hyperthermoplanets might only support extremophilic life.
The HEC uses the following terms to classify exoplanets in terms of mass, from least to greatest: asteroidan, mercurian, subterran, terran, superterran, neptunian, and jovian.
Alternative star systems
In determining the feasibility of extraterrestrial life, astronomers had long focused their attention on stars like the Sun. However, since planetary systems that resemble the Solar System are proving to be rare, they have begun to explore the possibility that life might form in systems very unlike the Sun's.
It is believed that F, G, K and M-type stars could host habitable exoplanets. About half of the stars similar in temperature to the Sun could have a rocky planet able to support liquid water on its surface, according to research using data from NASA's Kepler Space Telescope.
Binary systems
Typical estimates often suggest that 50% or more of all stellar systems are binary systems. This may be partly sample bias, as massive and bright stars tend to be in binaries and these are most easily observed and catalogued; a more precise analysis has suggested that the more common fainter stars are usually singular, and that up to two thirds of all stellar systems are therefore solitary.
The separation between stars in a binary may range from less than one astronomical unit (AU, the average Earth–Sun distance) to several hundred. In latter instances, the gravitational effects will be negligible on a planet orbiting an otherwise suitable star and habitability potential will not be disrupted unless the orbit is highly eccentric. However, where the separation is significantly less, a stable orbit may be impossible. If a planet's distance to its primary exceeds about one fifth of the closest approach of the other star, orbital stability is not guaranteed. Whether planets might form in binaries at all had long been unclear, given that gravitational forces might interfere with planet formation. Theoretical work by Alan Boss at the Carnegie Institution has shown that gas giants can form around stars in binary systems much as they do around solitary stars.
One study of Alpha Centauri, the nearest star system to the Sun, suggested that binaries need not be discounted in the search for habitable planets. Centauri A and B have an 11 AU distance at closest approach (23 AU mean), and both should have stable habitable zones. A study of long-term orbital stability for simulated planets within the system shows that planets within approximately three AU of either star may remain rather stable (i.e. the semi-major axis deviating by less than 5% during 32 000 binary periods). The continuous habitable zone (CHZ for 4.5 billion years) for Centauri A is conservatively estimated at 1.2 to 1.3 AU and Centauri B at 0.73 to 0.74—well within the stable region in both cases.
Red dwarf systems
M-type stars also considered possible hosts of habitable exoplanets, even those with flares such as Proxima b. Determining the habitability of red dwarf stars could help determine how common life in the universe might be, as red dwarfs make up between 70 and 90% of all the stars in the galaxy. However, it is important to bear in mind that flare stars could greatly reduce the habitability of exoplanets by eroding their atmosphere.
Size
Astronomers for many years ruled out red dwarfs as potential abodes for life. Their small size (from 0.08 to 0.45 solar masses) means that their nuclear reactions proceed exceptionally slowly, and they emit very little light (from 3% of that produced by the Sun to as little as 0.01%). Any planet in orbit around a red dwarf would have to huddle very close to its parent star to attain Earth-like surface temperatures; from 0.3 AU (just inside the orbit of Mercury) for a star like Lacaille 8760, to as little as 0.032 AU for a star like Proxima Centauri (such a world would have a year lasting just 6.3 days). At those distances, the star's gravity would cause tidal locking. One side of the planet would eternally face the star, while the other would always face away from it. The only ways in which potential life could avoid either an inferno or a deep freeze would be if the planet had an atmosphere thick enough to transfer the star's heat from the day side to the night side, or if there was a gas giant in the habitable zone, with a habitable moon, which would be locked to the planet instead of the star, allowing a more even distribution of radiation over the moon. It was long assumed that such a thick atmosphere would prevent sunlight from reaching the surface in the first place, preventing photosynthesis.
This pessimism has been tempered by research. Studies by Robert Haberle and Manoj Joshi of NASA's Ames Research Center in California have shown that a planet's atmosphere (assuming it included greenhouse gases CO2 and H2O) need only be , for the star's heat to be effectively carried to the night side. This is well within the levels required for photosynthesis, though water would still remain frozen on the dark side in some of their models. Martin Heath of Greenwich Community College, has shown that seawater, too, could be effectively circulated without freezing solid if the ocean basins were deep enough to allow free flow beneath the night side's ice cap. Further research—including a consideration of the amount of photosynthetically active radiation—suggested that tidally locked planets in red dwarf systems might at least be habitable for higher plants.
Other factors limiting habitability
Size is not the only factor in making red dwarfs potentially unsuitable for life, however. On a red dwarf planet, photosynthesis on the night side would be impossible, since it would never see the sun. On the day side, because the sun does not rise or set, areas in the shadows of mountains would remain so forever. Photosynthesis as we understand it would be complicated by the fact that a red dwarf produces most of its radiation in the infrared, and on the Earth the process depends on visible light. There are potential positives to this scenario. Numerous terrestrial ecosystems rely on chemosynthesis rather than photosynthesis, for instance, which would be possible in a red dwarf system. A static primary star position removes the need for plants to steer leaves toward the sun, deal with changing shade/sun patterns, or change from photosynthesis to stored energy during night. Because of the lack of a day-night cycle, including the weak light of morning and evening, far more energy would be available at a given radiation level.
Red dwarfs are far more variable and violent than their more stable, larger cousins. Often they are covered in starspots that can dim their emitted light by up to 40% for months at a time, while at other times they emit gigantic flares that can double their brightness in a matter of minutes. Such variation would be very damaging for life, as it would not only destroy any complex organic molecules that could possibly form biological precursors, but also because it would blow off sizeable portions of the planet's atmosphere.
For a planet around a red dwarf star to support life, it would require a rapidly rotating magnetic field to protect it from the flares. A tidally locked planet rotates only very slowly, and so cannot produce a geodynamo at its core. The violent flaring period of a red dwarf's life cycle is estimated to only last roughly the first 1.2 billion years of its existence. If a planet forms far away from a red dwarf so as to avoid tidal locking, and then migrates into the star's habitable zone after this turbulent initial period, it is possible that life may have a chance to develop. However, observations of the 7 to 12-billion year old Barnard's Star showcase that even old red dwarfs can have significant flare activity. Barnard's Star was long assumed to have little activity, but in 1998 astronomers observed an intense stellar flare, showing that it is a flare star.
Longevity and ubiquity
Red dwarfs have one advantage over other stars as abodes for life: far greater longevity. It took 4.5 billion years before humanity appeared on Earth, and life as we know it will see suitable conditions for 1 to 2.3 more. Red dwarfs, by contrast, could live for trillions of years because their nuclear reactions are far slower than those of larger stars, meaning that life would have longer to evolve and survive.
While the likelihood of finding a planet in the habitable zone around any specific red dwarf is slight, the total amount of habitable zone around all red dwarfs combined is equal to the total amount around Sun-like stars given their ubiquity. Furthermore, this total amount of habitable zone will last longer, because red dwarf stars live for hundreds of billions of years or even longer on the main sequence. However, combined with the above disadvantages, it is more likely that red dwarf stars would remain habitable longer to microbes, while the shorter-lived yellow dwarf stars, like the Sun, would remain habitable longer to animals.
Massive stars
Recent research suggests that very large stars, greater than ~100 solar masses, could have planetary systems consisting of hundreds of Mercury-sized planets within the habitable zone. Such systems could also contain brown dwarfs and low-mass stars (~0.1–0.3 solar masses). However the very short lifespans of stars of more than a few solar masses would scarcely allow time for a planet to cool, let alone the time needed for a stable biosphere to develop. Massive stars are thus eliminated as possible abodes for life.
However, a massive-star system could be a progenitor of life in another way – the supernova explosion of the massive star in the central part of the system. This supernova will disperse heavier elements throughout its vicinity, created during the phase when the massive star has moved off of the main sequence, and the systems of the potential low-mass stars (which are still on the main sequence) within the former massive-star system may be enriched with the relatively large supply of the heavy elements so close to a supernova explosion. However, this states nothing about what types of planets would form as a result of the supernova material, or what their habitability potential would be.
Neutron stars
Post-main sequence stars
Four classes of habitable planets based on water
In a review of the factors which are important for the evolution of habitable Earth-sized planets, Lammer et al. proposed a classification of four water-dependent habitat types:
Class I habitats are planetary bodies on which stellar and geophysical conditions allow liquid water to be available at the surface, along with sunlight, so that complex multicellular organisms may originate.
Class II habitats include bodies which initially enjoy Earth-like conditions, but do not keep their ability to sustain liquid water on their surface due to stellar or geophysical conditions. Mars, and possibly Venus are examples of this class where complex life forms may not develop.
Class III habitats are planetary bodies where liquid water oceans exist below the surface, where they can interact directly with a silicate-rich core.
Such a situation can be expected on water-rich planets located too far from their star to allow surface liquid water, but on which subsurface water is in liquid form because of the geothermal heat. Two examples of such an environment are Europa and Enceladus. In such worlds, not only is light not available as an energy source, but the organic material brought by meteorites (thought to have been necessary to start life in some scenarios) may not easily reach the liquid water. If a planet can only harbor life below its surface, the biosphere would not likely modify the whole planetary environment in an observable way, thus, detecting its presence on an exoplanet would be extremely difficult.
Class IV habitats have liquid water layers between two ice layers, or liquids above ice.
If the water layer is thick enough, water at its base will be in solid phase (ice polymorphs) because of the high pressure. Ganymede and Callisto are likely examples of this class. Their oceans are thought to be enclosed between thick ice layers. In such conditions, the emergence of even simple life forms may be very difficult because the necessary ingredients for life will likely be completely diluted.
The galactic neighborhood
Along with the characteristics of planets and their star systems, the wider galactic environment may also impact habitability. Scientists considered the possibility that particular areas of galaxies (galactic habitable zones) are better suited to life than others; the Solar System, in the Orion Arm, on the Milky Way galaxy's edge is considered to be in a life-favorable spot:
It is not in a globular cluster where immense star densities are inimical to life, given excessive radiation and gravitational disturbance. Globular clusters are also primarily composed of older, probably metal-poor, stars. Furthermore, in globular clusters, the great ages of the stars would mean a large amount of stellar evolution by the host or other nearby stars, which due to their proximity may cause extreme harm to life on any planets, provided that they can form.
It is not near an active gamma ray source.
It is not near the galactic center where once again star densities increase the likelihood of ionizing radiation (e.g., from magnetars and supernovae). The supermassive black holes at the centers of galaxies may also prove a danger to any nearby bodies.
The circular orbit of the Sun around the galactic center keeps it out of the way of the galaxy's spiral arms where intense radiation and gravitation may again lead to disruption.
Thus, relative isolation is ultimately what a life-bearing system needs. If the Sun were crowded amongst other systems, the chance of being fatally close to dangerous radiation sources would increase significantly. Further, close neighbors might disrupt the stability of various orbiting bodies such as Oort cloud and Kuiper belt objects, which can bring catastrophe if knocked into the inner Solar System.
While stellar crowding proves disadvantageous to habitability, so too does extreme isolation. A star as metal-rich as the Sun would probably not have formed in the very outermost regions of the Milky Way given a decline in the relative abundance of metals and a general lack of star formation. Thus, a "suburban" location, such as the Solar System enjoys, is preferable to a Galaxy's center or farthest reaches.
Other considerations
Alternative biochemistries
While most investigations of extraterrestrial life start with the assumption that advanced life-forms must have similar requirements for life as on Earth, the hypothesis of other types of biochemistry suggests the possibility of lifeforms evolving around a different metabolic mechanism. In Evolving the Alien, biologist Jack Cohen and mathematician Ian Stewart argue astrobiology, based on the Rare Earth hypothesis, is restrictive and unimaginative. They suggest that Earth-like planets may be very rare, but non-carbon-based complex life could possibly emerge in other environments. The most frequently mentioned alternative to carbon is silicon-based life, while ammonia and hydrocarbons are sometimes suggested as alternative solvents to water. The astrobiologist Dirk Schulze-Makuch and other scientists have proposed a Planet Habitability Index whose criteria include "potential for holding a liquid solvent" that is not necessarily restricted to water.
More speculative ideas have focused on bodies altogether different from Earth-like planets. Astronomer Frank Drake, a well-known proponent of the search for extraterrestrial life, imagined life on a neutron star: submicroscopic "nuclear molecules" combining to form creatures with a life cycle millions of times quicker than Earth life. Called "imaginative and tongue-in-cheek", the idea gave rise to science fiction depictions. Carl Sagan, another optimist with regards to extraterrestrial life, considered the possibility of organisms that are always airborne within the high atmosphere of Jupiter in a 1976 paper. Cohen and Stewart also envisioned life in both a solar environment and in the atmosphere of a gas giant.
"Good Jupiters"
"Good Jupiters" are gas giants, like the Solar System's Jupiter, that orbit their stars in circular orbits far enough away from the habitable zone not to disturb it but close enough to "protect" terrestrial planets in closer orbit in two critical ways. First, they help to stabilize the orbits, and thereby the climates of the inner planets. Second, they keep the inner stellar system relatively free of comets and asteroids that could cause devastating impacts. Jupiter orbits the Sun at about five times the distance between the Earth and the Sun. This is the rough distance we should expect to find good Jupiters elsewhere. Jupiter's "caretaker" role was dramatically illustrated in 1994 when Comet Shoemaker–Levy 9 impacted the giant.
However, the evidence is not quite so clear. Research has shown that Jupiter's role in determining the rate at which objects hit Earth is significantly more complicated than once thought.
The role of Jupiter in the early history of the Solar System is somewhat better established, and the source of significantly less debate. Early in the Solar System's history, Jupiter is accepted as having played an important role in the hydration of our planet: it increased the eccentricity of asteroid belt orbits and enabled many to cross Earth's orbit and supply the planet with important volatiles such as water and carbon dioxide. Before Earth reached half its present mass, icy bodies from the Jupiter–Saturn region and small bodies from the primordial asteroid belt supplied water to the Earth due to the gravitational scattering of Jupiter and, to a lesser extent, Saturn. Thus, while the gas giants are now helpful protectors, they were once suppliers of critical habitability material.
In contrast, Jupiter-sized bodies that orbit too close to the habitable zone but not in it (as in 47 Ursae Majoris), or have a highly elliptical orbit that crosses the habitable zone (like 16 Cygni B) make it very difficult for an independent Earth-like planet to exist in the system. See the discussion of a stable habitable zone above. However, during the process of migrating into a habitable zone, a Jupiter-size planet may capture a terrestrial planet as a moon. Even if such a planet is initially loosely bound and following a strongly inclined orbit, gravitational interactions with the star can stabilize the new moon into a close, circular orbit that is coplanar with the planet's orbit around the star.
Life's impact on habitability
A supplement to the factors that support life's emergence is the notion that life itself, once formed, becomes a habitability factor in its own right. An important Earth example was the production of molecular oxygen gas () by ancient cyanobacteria, and eventually photosynthesizing plants, leading to a radical change in the composition of Earth's atmosphere. This environmental change is called the Great Oxidation Event. This oxygen proved fundamental to the respiration of later animal species. The Gaia hypothesis, a scientific model of the geo-biosphere pioneered by James Lovelock in 1975, argues that life as a whole fosters and maintains suitable conditions for itself by helping to create a planetary environment suitable for its continuity. Similarly, David Grinspoon has suggested a "living worlds hypothesis" in which our understanding of what constitutes habitability cannot be separated from life already extant on a planet. Planets that are geologically and meteorologically alive are much more likely to be biologically alive as well and "a planet and its life will co-evolve." This is the basis of Earth system science.
The role of chance
In 2020, a computer simulation of the evolution of planetary climates over 3 billion years suggested that feedback is a necessary but insufficient condition for preventing planets from ever becoming too hot or cold for life. Chance also plays a crucial role. Related considerations include yet unknown factors influencing the thermal habitability of planets such as "feedback mechanism (or mechanisms) that prevents the climate ever wandering to fatal temperatures".
| Physical sciences | Planetary science | Astronomy |
2592962 | https://en.wikipedia.org/wiki/Octant%20%28instrument%29 | Octant (instrument) | The octant, also called a reflecting quadrant, is a reflecting instrument used in navigation.
Etymology
The name octant derives from the Latin octans meaning eighth part of a circle, because the instrument's arc is one eighth of a circle.
Reflecting quadrant derives from the instrument using mirrors to reflect the path of light to the observer and, in doing so, doubles the angle measured. This allows the instrument to use a one-eighth of a turn to measure a quarter-turn or quadrant.
Origin of the octant
Newton's reflecting quadrant
Isaac Newton's reflecting quadrant was invented around 1699. A detailed description of the instrument was given to Edmond Halley, but the description was not published until after Halley's death in 1742. It is not known why Halley did not publish the information during his life, as this prevented Newton from getting the credit for the invention that is generally given to John Hadley and Thomas Godfrey.
One copy of this instrument was constructed by Thomas Heath (instrument maker) and may have been shown in Heath's shop window prior to its being published by the Royal Society in 1742.
Newton's instrument used two mirrors, but they were used in an arrangement somewhat different from the two mirrors found in modern octants and sextants. The diagram on the right shows the configuration of the instrument.
The 45° arc of the instrument (PQ), was graduated with 90 divisions of a half-degree each. Each such division was subdivided into 60 parts and each part further divided into sixths. This results in the arc being marked in degrees, minutes and sixths of a minute (10 seconds). Thus the instrument could have readings interpolated to 5 seconds of arc. This fineness of graduation is only possible due to the large size of the instrument - the sighting telescope alone was three to four feet long.
A sighting telescope (AB), three or four feet long, was mounted along one side of the instrument. A horizon mirror was fixed at a 45° angle in front of the telescope's objective lens (G). This mirror was small enough to allow the observer to see the image in the mirror on one side and to see directly ahead on the other. The index arm (CD) held an index mirror (H), also at 45° to the edge of the index arm. The reflective sides of the two mirrors nominally faced each other, so that the image seen in the first mirror is that reflected from the second.
With the two mirrors parallel, the index reads 0°. The view through the telescope sees directly ahead on one side and the view from the mirror G sees the same image reflected from mirror H (see detail drawing to the right). When the index arm is moved from zero to a large value, the index mirror reflects an image that is in a direction away from the direct line of sight. As the index arm movement increases, the line of sight for the index mirror moves toward S (to the right in the detail image). This shows a slight deficiency with this mirror arrangement. The horizon mirror will block the view of the index mirror at angles approaching 90°.
The length of the sighting telescope seems remarkable, given the small size of the telescopes on modern instruments. This was likely Newton's choice of a way to reduce chromatic aberrations. Short–focal length telescopes, prior to the development of achromatic lenses, produced an objectionable degree of aberration, so much so that it could affect the perception of a star's position. Long focal lengths were the solution, and this telescope would likely have had both a long–focal length objective lens and a long–focal length eyepiece. This would decrease aberrations without excessive magnification.
The inventors of the octant
Two men independently developed the octant around 1730: John Hadley (1682–1744), an English mathematician, and Thomas Godfrey (1704–1749), a glazier in Philadelphia. While both have a legitimate and equal claim to the invention, Hadley generally gets the greater share of the credit. This reflects the central role that London and the Royal Society played in the history of scientific instruments in the eighteenth century.
Two others who created octants during this period were Caleb Smith, an English insurance broker with a strong interest in astronomy (in 1734), and Jean-Paul Fouchy, a mathematics professor and astronomer in France (in 1732).
Hadley's versions
Hadley produced two versions of the reflecting quadrant. Only the second is well known and is the familiar octant.
Hadley's reflecting quadrant
Hadley's first reflecting quadrant was a simple device with a frame spanning a 45° arc. In the image at the right, from Hadley's article in the Philosophical Transactions of the Royal Society, you can see the nature of his design. A small sighting telescope was mounted on the frame along one side. One large index mirror was mounted at the point of rotation of the index arm. A second, smaller horizon mirror was mounted on the frame in the line of sight of the telescope. The horizon mirror allows the observer to see the image of the index mirror in one half of the view and to see a distant object in the other half. A shade was mounted at the vertex of the instrument to allow one to observe a bright object. The shade pivots to allow it to move out of the way for stellar observations.
Observing through the telescope, the navigator would sight one object directly ahead. The second object would be seen by reflection in the horizon mirror. The light in the horizon mirror is reflected from the index mirror. By moving the index arm, the index mirror can be made to reveal any object up to 90° from the direct line of sight. When both objects are in the same view, aligning them together allows the navigator to measure the angular distance between them.
Very few of the original reflecting quadrant designs were ever produced. One, constructed by Baradelle, is in the collection of the Musée de la Marine, Paris.
Hadley's octant
Hadley's second design had the form familiar to modern navigators. The image to the right, also taken from his Royal Society publication, shows the details.
He placed an index mirror on the index arm. Two horizon mirrors were provided. The upper mirror, in the line of the sighting telescope, was small enough to allow the telescope to see directly ahead as well as seeing the reflected view. The reflected view was that of the light from the index mirror. As in the previous instrument, the arrangement of the mirrors allowed the observer to simultaneously see an object straight ahead and to see one reflected in the index mirror to the horizon mirror and then into the telescope. Moving the index arm allowed the navigator to see any object within 90° of the direct view.
The significant difference with this design was that the mirrors allowed the instrument to be held vertically rather than horizontally and it provided more room for configuring the mirrors without suffering from mutual interference.
The second horizon mirror was an interesting innovation. The telescope was removable. It could be remounted so that the telescope viewed the second horizon mirror from the opposite side of the frame. By mounting the two horizon mirrors at right angles to each other and permitting the movement of the telescope, the navigator could measure angles from 0 to 90° with one horizon mirror and from 90° to 180° with the other. This made the instrument very versatile. For unknown reasons, this feature was not implemented on octants in general use.
Comparing this instrument to the photo of a typical octant at the top of the article, one can see that the only significant differences in the more modern design are:
The location of the horizon mirror and telescope or sighting pinnula is lower.
The internal bracing of the frame is more central and robust.
The position of the shades for the index mirror is in the path between the index and horizon mirrors rather than at the top of the instrument.
Multiple shades are used to allow for different levels of shading.
Separate shades are provided on the horizon mirror for sighting a low sun position with a very bright horizon.
The second horizon mirror and accompanying alidade is not provided.
Smith's Astroscope
Caleb Smith, an English insurance broker with a strong interest in astronomy, had created an octant in 1734. He called it an Astroscope or Sea-Quadrant. His used a fixed prism in addition to an index mirror to provide reflective elements. Prisms provide advantages over mirrors in an era when polished speculum metal mirrors were inferior and both the silvering of a mirror and the production of glass with flat, parallel surfaces was difficult.
In the drawing to the right, the horizon element (B) could be a mirror or a prism. On the index arm, the index mirror (A) rotated with the arm. A sighting telescope was mounted on the frame (C). The index did not use a vernier or other device at the scale (D). Smith called the instrument's index arm a label, in the manner of Elton for his mariner's quadrant.
Various design elements of Smith's instrument made it inferior to Hadley's octant and it was not used significantly. For example, one problem with the Astroscope was that angle of the observer's line of sight. By looking down, he had greater difficulty in observing than an orientation with his head in a normal orientation.
Advantages of the octant
The octant provided a number of advantages over previous instruments.
The sight was easy to align because the horizon and the star seem to move together as the ship pitched and rolled. This also created a situation where the error in observation was less dependent on the observer, as they could directly see both objects at once.
With the use of the manufacturing techniques available in the 18th century, the instruments were capable of reading very accurately. The size of the instruments was reduced with no loss of accuracy. An octant could be half the size of a Davis quadrant with no increase in error.
Using shades over the light paths, one could observe the sun directly, while moving the shades out of the light path allowed the navigator to observe faint stars. This made the instrument usable both night and day.
By 1780, the octant and sextant had almost completely displaced all previous navigational instruments.
Production of the octant
Early octants were constructed primarily in wood, with later versions incorporating ivory and brass components. The earliest mirrors were polished metal, since the technology to produce silvered glass mirrors with flat, parallel surfaces was limited. As glass polishing techniques improved, glass mirrors began to be provided. These used coatings of mercury-containing tin amalgam; coatings of silver or aluminum were not available until the 19th century. The poor optical quality of the early polished speculum metal mirrors meant that telescopic sights were not practical. For that reason, most early octants employed a simple naked-eye sighting pinnula instead.
Early octants retained some of the features common to backstaves, such as transversals on the scale. However, as engraved, they showed the instrument to have an apparent accuracy of only two minutes of arc while the backstaff appeared to be accurate to one minute. The use of the vernier scale allowed the scale to be read to one minute, so improved the marketability of the instrument. This and the ease in making verniers compared to transversals, lead to adoption of the vernier on octants produced later in the 18th century.
Octants were produced in large numbers. In wood and ivory, their relatively low price compared to an all-brass sextant made them a popular instrument. The design was standardized with many manufacturers using the identical frame style and components. Different shops could make different components, with woodworkers specializing in frames and others in the brass components. For example, Spencer, Browning and Rust, a manufacturer of scientific instruments in England from 1787 to 1840 (operating as Spencer, Browning and Co. after 1840) used a Ramsden dividing engine to produce graduated scales in ivory. These were widely used by others and the SBR initials could be found on octants from many other manufacturers.
Examples of these very similar octants are in the photos in this article. The image at the top is essentially the same instrument as the one in the detail photos. However, they are from two different instrument makers - the upper is labelled Crichton - London, Sold by J Berry Aberdeen while the detail images are of an instrument from Spencer, Browning & Co. London. The only obvious difference is the presence of horizon shades on the Crichton octant that are not on the other.
These octants were available with many options. A basic octant with graduations directly on the wood frame were least expensive. These dispensed with a telescopic sight, using a single- or double-holed sighting pinnula instead. Ivory scales would increase the price, as would the use of a brass index arm or a vernier.
Demise of the octant
In 1767 the first edition of The Nautical Almanac tabulated lunar distances, enabling navigators to find the current time from the angle between the Sun and the Moon. This angle is sometimes larger than 90°, and thus not possible to measure with an octant. For that reason, Admiral John Campbell, who conducted shipboard experiments with the lunar distance method, suggested a larger instrument and the sextant was developed.
From that time onward, the sextant was the instrument that experienced significant development and improvements and was the instrument of choice for naval navigators. The octant continued to be produced well into the 19th century, though it was generally a less accurate and less expensive instrument. The lower price of the octant, including versions without telescope, made it a practical instrument for ships in the merchant and fishing fleets.
One common practice among navigators up to the late nineteenth century was to use both a sextant and an octant. The sextant was used with great care and only for lunars, while the octant was used for routine meridional altitude measurements of the Sun every day. This protected the very accurate and pricier sextant, while using the more affordable octant where it performs well.
Bubble octant
From the early 1930s through the end of the 1950s, several types of civilian and military bubble octant instruments were produced for use aboard aircraft. All were fitted with an artificial horizon in the form of a bubble, which was centered to align the horizon for a navigator flying thousands of feet above the Earth; some had recording features.
Use and adjustment
Use and adjustment of the octant is essentially identical to the navigator's sextant.
Other reflecting instruments
Hadley's was not the first reflecting quadrant. Robert Hooke invented a reflecting quadrant in 1684 and had written about the concept as early as 1666. Hooke's was a single-reflecting instrument. Other octants were developed by Jean-Paul Fouchy and Caleb Smith in the early 1730s, however, these did not become significant in the history of navigation instruments.
| Technology | Navigation | null |
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