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475317 | https://en.wikipedia.org/wiki/Ocellated%20turkey | Ocellated turkey | The ocellated turkey (Meleagris ocellata) is a species of turkey residing primarily in the Yucatán Peninsula, Mexico, as well as in parts of Belize and Guatemala. A relative of the North American wild turkey (Meleagris gallopavo), it was sometimes previously considered in a genus of its own (Agriocharis), but the differences between the two turkeys are currently considered too small to justify generic segregation. It is a relatively large bird, at around long and an average weight of in females and in males.
The ocellated turkey lives only in a range in the Yucatán Peninsula in Mexico—which includes all or part of the states of Quintana Roo, Campeche, Yucatán, Tabasco, and Chiapas—as well as the northern and western parts of Belize and northern Guatemala.
The ocellated turkey was considered endangered by Mexican authorities as recently as 2002 and has been considered Near Threatened by the IUCN since 2004.
Ecology
The species is believed to have experienced a decline in response to land use changes and higher than sustainable harvest by migrant workers and subsistence hunters living in the Yucatán Peninsula region of Central America. A study conducted in the year 2011 indicated that the ocellated turkey made up a substantial amount of the diets of four prominent ethnic groups of the Yucatán Peninsula.
Description
The body feathers of both sexes are a mixture of bronze and green iridescent color. Although females can be duller with more green, the breast feathers do not generally differ and cannot be used to determine sex. Neither sex possesses the "beard" typically found in wild turkeys. Tail feathers of both sexes are bluish-grey with an eye-shaped, blue-bronze spot near the end with a bright gold tip. The spots, or ocelli (located on the tail), for which the ocellated turkey is named, have been likened to the patterning typically found on peafowl. The upper, major secondary wing coverts are rich iridescent copper. The primary and secondary wing feathers have similar barring to that of North American turkeys, but the secondaries have more white, especially around the edges.
Both sexes have blue heads with some orange or red nodules, which are more pronounced on males. The males also have a fleshy blue crown covered with nodules, similar to those on the neck, behind the snood. During breeding season this crown swells up and becomes brighter and more pronounced in its yellow-orange color. The eye is surrounded by a ring of bright red skin, which is most visible on males during breeding season. The legs are deep red and are shorter and thinner than on North American turkeys. Males over one year old have spurs on the legs that average , with lengths of over being recorded. These spurs are much longer and thinner than on North American turkeys.
Ocellated turkeys are much smaller than any of the subspecies of North American wild turkey, with adult hens weighing about before laying eggs and 3 kg (6–7 pounds) the rest of the year, and adult males weighing about during breeding season.
Vocalizations of the hen ocellated turkeys are similar to those of their northern relatives, however the male vocalization known as a "gobble" is quite different in comparison. The gobble begins with several low frequency "thumps", much like the sound of a small gasoline motor starting. As the tempo of thumps increases, the typical gobble is produced. Branton and Berryhill (2007) have observed that the male ocellated turkey does not gobble per se like the wild turkey. Rather, his song is distinct and includes some six to seven bongo-like bass tones which quicken in both cadence and volume until a crescendo is reached whereupon the bird's head is fully erect while he issues forth a rather high-pitched but melodious series of chops. The ocellated turkey will typically begin his singing 20 to 25 minutes before sunrise—similar to the wild turkey.
Behavior
Turkeys spend most of the time on the ground and during the day, they often prefer running to escape danger rather than taking off, though they can fly swiftly and powerfully for short distances when necessary, as typical of gamebirds. Roosting is usually high in trees away from night-hunting predators, such as jaguars, and usually in a family group.
The ocellated turkey is a generalist in terms of its feeding habits. They are known to feed on a wide variety of forage including but not limited to insects, such as beetles, moths, and leafcutter ants, grass seeds, nuts, and leaves. The feeding rates of male ocellated turkeys have been observed to be significantly higher during January before the breeding season begins than when the breeding season is in full swing. Prior to the breeding season, adult male turkeys have been observed in flocks usually no larger than three mature birds, whereas flocks of eight or more birds consisted of yearling turkeys and hens.
The breeding season for the ocellated turkey begins in early February when the first gobbles are heard. The breeding season peaks in March and comes to an end by the end of April. Male ocellated turkeys engage in an elaborate, spirited display to attract females. Ocellated turkeys use their tail fans similar to North American turkeys; however, there are several distinct differences between the display of the ocellated and their North American cousins. Male turkeys begin the mating dance by tapping their feet against the ground in rapid succession. Next, the male birds move their tail feathers from side to side while quickly vibrating their wings and dragging the tips of them against the ground. As the male does this dance, he moves around the female making sure the dorsal surface of the tail feathers are constantly in view of the female.
Female ocellated turkeys lay 8–15 eggs in a well concealed nest on the ground. The poults hatch in May through July after a 28-day incubation period. The poults are covered in a reddish brown juvenile plumage which allows them to blend into their surroundings to hide from predators. The young are precocial and able to leave the nest after one night. They then follow their mother until they reach young adulthood when they begin to range though often re-grouping to roost.
In culture
The Mayans (who call it kuuts or yuum kuuts) have a story where the turkey's feathers were given to him by the nightjar who vouched for him to be anointed king of the birds by the Great Ancestor (Nohochacyum) in return for a reward. As the turkey completely forgot what he owed the nightjar, the latter complained to Nohochacyum who punished the former by changing his once-melodious voice to guttural thumps. It is associated with water, rain (especially through the deity Chaak) and fecundity; tribes like the Chʼortiʼ once offered its blood across their fields in hopes of a good harvest. Its carcasses were also offered to use by Mayan chiefs, even found in their tombs.
| Biology and health sciences | Galliformes | Animals |
475598 | https://en.wikipedia.org/wiki/Ecological%20fallacy | Ecological fallacy | An ecological fallacy (also ecological inference fallacy or population fallacy) is a formal fallacy in the interpretation of statistical data that occurs when inferences about the nature of individuals are deduced from inferences about the group to which those individuals belong. "Ecological fallacy" is a term that is sometimes used to describe the fallacy of division, which is not a statistical fallacy. The four common statistical ecological fallacies are: confusion between ecological correlations and individual correlations, confusion between group average and total average, Simpson's paradox, and confusion between higher average and higher likelihood. From a statistical point of view, these ideas can be unified by specifying proper statistical models to make formal inferences, using aggregate data to make unobserved relationships in individual level data.
Examples
Mean and median
An example of ecological fallacy is the assumption that a population mean has a simple interpretation when considering likelihoods for an individual.
For instance, if the mean score of a group is larger than zero, this does not imply that a random individual of that group is more likely to have a positive score than a negative one (as long as there are more negative scores than positive scores an individual is more likely to have a negative score). Similarly, if a particular group of people is measured to have a lower mean IQ than the general population, it is an error to conclude that a randomly-selected member of the group is more likely than not to have a lower IQ than the mean IQ of the general population; it is also not necessarily the case that a randomly selected member of the group is more likely than not to have a lower IQ than a randomly-selected member of the general population. Mathematically, this comes from the fact that a distribution can have a positive mean but a negative median. This property is linked to the skewness of the distribution.
Consider the following numerical example:
Group A: 80% of people got 40 points and 20% of them got 95 points. The mean score is 51 points.
Group B: 50% of people got 45 points and 50% got 55 points. The mean score is 50 points.
If we pick two people at random from A and B, there are 4 possible outcomes:
A – 40, B – 45 (B wins, 40% probability – 0.8 × 0.5)
A – 40, B – 55 (B wins, 40% probability – 0.8 × 0.5)
A – 95, B – 45 (A wins, 10% probability – 0.2 × 0.5)
A – 95, B – 55 (A wins, 10% probability – 0.2 × 0.5)
Although Group A has a higher mean score, 80% of the time a random individual of A will score lower than a random individual of B.
Individual and aggregate correlations
Research dating back to Émile Durkheim suggests that predominantly Protestant localities have higher suicide rates than predominantly Catholic localities. According to Freedman, the idea that Durkheim's findings link, at an individual level, a person's religion to their suicide risk is an example of the ecological fallacy. A group-level relationship does not automatically characterize the relationship at the level of the individual.
Similarly, even if at the individual level, wealth is positively correlated to tendency to vote Republican in the United States, we observe that wealthier states tend to vote Democratic. For example, in the 2004 United States presidential election, the Republican candidate, George W. Bush, won the fifteen poorest states, and the Democratic candidate, John Kerry, won 9 of the 11 wealthiest states in the Electoral College. Yet 62% of voters with annual incomes over $200,000 voted for Bush, but only 36% of voters with annual incomes of $15,000 or less voted for Bush.
Aggregate-level correlation will differ from individual-level correlation if voting preferences are affected by the total wealth of the state even after controlling for individual wealth. The true driving factor in voting preference could be self-perceived relative wealth; perhaps those who see themselves as better off than their neighbours are more likely to vote Republican. In this case, an individual would be more likely to vote Republican if they became wealthier, but they would be more likely to vote for a Democrat if their neighbor's wealth increased (resulting in a wealthier state).
However, the observed difference in voting habits based on state- and individual-level wealth could also be explained by the common confusion between higher averages and higher likelihoods as discussed above. States may not be wealthier because they contain more wealthy people (i.e., more people with annual incomes over $200,000), but rather because they contain a small number of super-rich individuals; the ecological fallacy then results from incorrectly assuming that individuals in wealthier states are more likely to be wealthy.
Many examples of ecological fallacies can be found in studies of social networks, which often combine analysis and implications from different levels. This has been illustrated in an academic paper on networks of farmers in Sumatra.
Robinson's paradox
A 1950 paper by William S. Robinson computed the illiteracy rate and the proportion of the population born outside the US for each state and for the District of Columbia, as of the 1930 census. He showed that these two figures were associated with a negative correlation of −0.53; in other words, the greater the proportion of immigrants in a state, the lower its average illiteracy (or, equivalently, the higher its average literacy). However, when individuals are considered, the correlation between illiteracy and nativity was +0.12 (immigrants were on average more illiterate than native citizens). Robinson showed that the negative correlation at the level of state populations was because immigrants tended to settle in states where the native population was more literate. He cautioned against deducing conclusions about individuals on the basis of population-level, or "ecological" data. In 2011, it was found that Robinson's calculations of the ecological correlations are based on the wrong state level data. The correlation of −0.53 mentioned above is in fact −0.46. Robinson's paper was seminal, but the term 'ecological fallacy' was not coined until 1958 by Selvin.
Formal problem
The correlation of aggregate quantities (or ecological correlation) is not equal to the correlation of individual quantities. Denote by Xi, Yi two quantities at the individual level. The formula for the covariance of the aggregate quantities in groups of size N is
The covariance of two aggregated variables depends not only on the covariance of two variables within the same individuals but also on covariances of the variables between different individuals. In other words, correlation of aggregate variables take into account cross sectional effects which are not relevant at the individual level.
The problem for correlations entails naturally a problem for regressions on aggregate variables: the correlation fallacy is therefore an important issue for a researcher who wants to measure causal impacts. Start with a regression model where the outcome is impacted by
The regression model at the aggregate level is obtained by summing the individual equations:
Nothing prevents the regressors and the errors from being correlated at the aggregate level. Therefore, generally, running a regression on aggregate data does not estimate the same model than running a regression with individual data.
The aggregate model is correct if and only if
This means that, controlling for , does not determine .
Choosing between aggregate and individual inference
There is nothing wrong in running regressions on aggregate data if one is interested in the aggregate model. For instance, for the governor of a state, it is correct to run regressions between police force on crime rate at the state level if one is interested in the policy implication of a rise in police force. However, an ecological fallacy would happen if a city council deduces the impact of an increase in police force in the crime rate at the city level from the correlation at the state level.
Choosing to run aggregate or individual regressions to understand aggregate impacts on some policy depends on the following trade-off: aggregate regressions lose individual-level data but individual regressions add strong modeling assumptions. Some researchers suggest that the ecological correlation gives a better picture of the outcome of public policy actions, thus they recommend the ecological correlation over the individual level correlation for this purpose (Lubinski & Humphreys, 1996). Other researchers disagree, especially when the relationships among the levels are not clearly modeled. To prevent ecological fallacy, researchers with no individual data can model first what is occurring at the individual level, then model how the individual and group levels are related, and finally examine whether anything occurring at the group level adds to the understanding of the relationship. For instance, in evaluating the impact of state policies, it is helpful to know that policy impacts vary less among the states than do the policies themselves, suggesting that the policy differences are not well translated into results, despite high ecological correlations (Rose, 1973).
Group and total averages
Ecological fallacy can also refer to the following fallacy: the average for a group is approximated by the average in the total population divided by the group size. Suppose one knows the number of Protestants and the suicide rate in the USA, but one does not have data linking religion and suicide at the individual level. If one is interested in the suicide rate of Protestants, it is a mistake to estimate it by the total suicide rate divided by the number of Protestants.
Formally, denote the mean of the group, we generally have:
However, the law of total probability gives
As we know that is between 0 and 1, this equation gives a bound for .
Simpson's paradox
A striking ecological fallacy is Simpson's paradox: the fact that when comparing two populations divided into groups, the average of some variable in the first population can be higher in every group and yet lower in the total population. Formally, when each value of Z refers to a different group and X refers to some treatment, it can happen that
When does not depend on , the Simpson's paradox is exactly the omitted variable bias for the regression of Y on X where the regressor is a dummy variable and the omitted variable is a categorical variable defining groups for each value it takes. The application is striking because the bias is high enough that parameters have opposite signs.
Legal applications
The ecological fallacy was discussed in a court challenge to the 2004 Washington gubernatorial election in which a number of illegal voters were identified, after the election; their votes were unknown, because the vote was by secret ballot. The challengers argued that illegal votes cast in the election would have followed the voting patterns of the precincts in which they had been cast, and thus adjustments should be made accordingly. An expert witness said this approach was like trying to figure out Ichiro Suzuki's batting average by looking at the batting average of the entire Seattle Mariners team, since the illegal votes were cast by an unrepresentative sample of each precinct's voters, and might be as different from the average voter in the precinct as Ichiro was from the rest of his team. The judge determined that the challengers' argument was an ecological fallacy and rejected it.
| Mathematics | Statistics | null |
475682 | https://en.wikipedia.org/wiki/Clostridium%20perfringens | Clostridium perfringens | Clostridium perfringens (formerly known as C. welchii, or Bacillus welchii) is a Gram-positive, bacillus (rod-shaped), anaerobic, spore-forming pathogenic bacterium of the genus Clostridium. C. perfringens is ever-present in nature and can be found as a normal component of decaying vegetation, marine sediment, the intestinal tract of humans and other vertebrates, insects, and soil. It has the shortest reported generation time of any organism at 6.3 minutes in thioglycolate medium.
Clostridium perfringens is one of the most common causes of food poisoning in the United States, alongside norovirus, Salmonella, Campylobacter, and Staphylococcus aureus. However, it can sometimes be ingested and cause no harm.
Infections induced by C. perfringens are associated with tissue necrosis, bacteremia, emphysematous cholecystitis, and gas gangrene, which is also known as clostridial myonecrosis. The specific name, perfringens, is derived from the Latin (meaning "through") and ("burst"), referring to the disruption of tissue that occurs during gas gangrene. Gas gangrene is caused by alpha toxin, or α-toxin, that embeds itself into the plasma membrane of cells and disrupts normal cellular function by altering membrane structure. Research suggests that C. perfringens is capable of engaging in polymicrobial anaerobic infections. It is commonly encountered in infections as a component of the normal flora. In this case, its role in disease is minor.
C. perfringens toxins are a result of horizontal gene transfer of a neighboring cell's plasmids. Shifts in genomic make-up are common for this species of bacterium and contribute to novel pathogenesis. Major toxins are expressed differently in certain populations of C. perfringens; these populations are organized into strains based on their expressed toxins. This especially impacts the food industry, as controlling this microbe is important for preventing foodborne illness. Novel findings in C. perfringens hyper-motility, which was provisionally thought as non-motile, have been discovered as well. Findings in metabolic processes reveal more information concerning C. perfringens pathogenic nature.
Genome
Clostridium perfringens has a stable G+C content around 27 to 28 percent and average genome size of 3.5 Mb. Genomes of 56 C. perfringens strains have since been made available on the NCBI genomes database for the scientific research community. Genomic research has revealed surprisingly high diversity in C. perfringens pangenome, with only 12.6 percent core genes, identified as the most divergent Gram-positive bacteria reported. Nevertheless, 16S rRNA regions in between C. perfringens strains are found to be highly conserved (sequence identity >99.1%).
The Clostridium perfringens enterotoxin (CPE)–producing strain has been identified to be a small portion of the overall C. perfringens population (~1-5%) through genomic testing. Advances in genetic information surrounding strain A CPE C. perfringens has allowed techniques such as microbial source tracking (MST) to identify food contamination sources. The CPE gene has been found within chromosomal DNA as well as plasmid DNA. Plasmid DNA has been shown to play and integral role in cell pathogenesis and encodes for major toxins, including CPE.
C. perfringens has been shown to carry plasmid-containing genes for antibiotic resistance. The pCW3 plasmid is the primary conjugation plasmid responsible for creating antibiotic resistance in C. perfringens. Furthermore, the pCW3 plasmid also encodes for multiple toxins found in pathogenic strains of C. perfringens. Antibiotic resistance genes observed thus far include tetracycline resistance, efflux protein, and aminoglycoside resistance.
Within industrial contexts, such as food production, sequencing genomes for pathogenic strains of C. perfringens has become an expanding field of research. Poultry production is impacted directly from this trend as antibiotic-resistant strains of C. perfringens are becoming more common. By performing a meta-genome analysis, researches are capable to identify novel strains of pathogenic bacterium, such as C. perfringens B20.
Motility
Clostridium perfringens is provisionally identified as non-motile. They lack flagella; however, recent research suggests gliding as a form of motility.
Hyper-motile variations
In agar plate cultures bacteria with hypermotile variations like SM101 frequently appear around the borders of the colonies. They create long thin filaments that enable them to move quickly, much like bacteria with flagella, according to video imaging of their gliding motion. The causes of the hypermotile phenotype and its immediate descendants were found using genome sequencing. The hypermotile offspring of strains SM101 and SM102, SM124 and SM127, respectively, had 10 and 6 nucleotide polymorphisms (SNPs) in comparison to their parent strains. The hypermotile strains have the common trait of gene mutations related to cell division.
Regulation of gliding motility: The CpAL/VirSR system
Some strains of C. perfringens cause various diseases like gas gangrene and myonecrosis. Toxins produced that are required for myonecrosis is regulated by the C. perfringens Agr-like (CpAl) system through the VirSR two-component system. The CpAL/VirSR system is a quorum sensing system encoded by other pathogenic clostridia. Myonecrosis starts at the infection site and involves bacteria migrating deeper via gliding motility. Researchers investigated if the CpAL/VirSR system regulates gliding motility. The study demonstrated that the CpAL/VirSR regulates C. perfringens gliding motility. Additionally, the study suggests that gliding bacteria in myonecrosis have increased transcription of toxin genes.
Transformation
There are two methods of genetic manipulation via experimentation that have been shown to cause genetic transformation in C. perfringens.
Protoplast transformation
The first report of transformation in C. perfringens involved polyethylene glycol-mediated transformation of protoplasts. The transformation procedure involved the addition of the plasmid DNA to the protoplasts in the presence of high concentrations of polyethylene glycol. During the first protoplast transformation experiment, L-phase variants of C. perfringens were generated by penicillin treatment in the presence 0.4m sucrose. After the transformation procedure was completed, all of the transformed cells were still in the form of L-phase variants. Reversion to vegetative cells was not obtained, but it was observed that autoplasts (protoplasts derived from autolysis) were able to be regenerated to produce rods with cell walls and could be transformed with C. perfringens plasmid DNA.
Electroporation
Electroporation involves the application of a high-voltage electric field to vegetative bacteria cells for a very short period. This technique resulted in major advances in genetic transformation of C. perfringens, due to the bacteria often displaying itself as a vegetative cell or as dormant spores in food. The electric pulse creates pores in the bacterial cell membrane and allows the passive influx of DNA molecules.
Metabolic processes
C. perfringens is an aerotolerant anaerobe bacterium that lives in a variety of environments including soil and human intestinal tract. C. perfringens is incapable of synthesizing multiple amino acids due to the lack of genes required for biosynthesis. Instead, the bacterium produces enzymes and toxins to break down host cells and import nutrients from the degrading cell.
C. perfringens has a complete set of enzymes for glycolysis and glycogen metabolism. In the fermentation pathway, pyruvate is converted into acetyl-CoA by pyruvate-ferredoxin oxidoreductase, producing CO2 gas and reduced ferredoxin. Electrons from the reduced ferredoxin are transferred to protons by hydrogenase, resulting in the formation of hydrogen molecules (H2) that are released from the cell along with CO2. Pyruvate is also converted to lactate by lactate dehydrogenase, whereas acetyl-CoA is converted into ethanol, acetate, and butyrate through various enzymatic reactions, completing the anaerobic glycolysis that serves as a potential main energy source for C. perfringens. C. perfringens utilizes a variety of sugars such as fructose, galactose, glycogen, lactose, maltose, mannose, raffinose, starch, and sucrose, and various genes for glycolytic enzymes. The amino acids of these various enzymes and sugar molecules are converted to propionate through propionyl-CoA, which results in energy production.
Virulence
Membrane-damaging enzymes, pore-forming toxins, intracellular toxins, and hydrolytic enzymes are the functional categories into which C. perfringens' virulence factors may be divided. These virulence factor-encoding genes can be found on chromosomes and large plasmids.
Carbohydrate-active enzymes
The human gastrointestinal tract is lined with intestinal mucosa that secrete mucus and act as a defense mechanism against pathogens, toxins, and harmful substances. Mucus is made up of mucins containing several O-linked glycan glycoproteins that recognizes and forms a barrier around microbes, preventing them from attaching to endothelial cells and infecting them. C. perfringens can secrete different carbohydrate-active enzymes (CAZymes) that aid in degrading mucins and other O-glycans within the intestinal mucosa. These enzymes include: Sialidases, Hexosaminidases, Galactosidases, and Fucosidases belonging to various glycoside hydrolase families.
Sialidase
Sialidases, also called neuraminidases, function to breakdown mucin by hydrolyzing the terminal sialic acid residues located within the protein through the process of desialylation. C. perfringens has three sialidases belonging to glycoside hydrolase family 33 (GH33): NanH, NanI, and NanJ. All strains of C. perfringens encode for at least one of these enzymes.
C. perfringens can secrete NanI and NanJ through secretion signal peptides located on each protein. Research suggests that NanH operates in the cytoplasm of C. perfringens, as it does not contain a secretion signal peptide. NanH contains only a catalytic domain, whereas NanI and NanJ contain a catalytic domain and additional carbohydrate-binding modules (CBMs) to aid in catalytic activity. Located on their N-terminals, NanI contains CBM40, whereas NanJ contains both CBM40 and CBM32. Based on studies analyzing the three-dimensional structure of NanI, its active site has a pocket-like orientation that aids in the removal of sialic acid residues from sialomucins in the intestinal mucosa.
Hexosaminidase
The mucus layer consists of intestinal mucin glycans, glycolipids, and glycoproteins that contain hexosamines, such as N-acetylglucosamine (GlcNAc) and N-acetylgalactosamine (GalNAc). C. perfringens encodes for eight hexosaminidases that break down hexosamines in the mucus. These hexosaminidases belong to four glycoside hydrolase families: GH36, GH84, GH89, and GH123.
C. perfringens encodes for AagA (CpGH36A) and CpGH36B in glycoside hydrolase family 36 (GH36): AagA removes GalNAc from O-glycans, and CpGH36B is expected to have a similar structure to AagA, but specificities on its function are unknown. NagH, NagI, NagJ, and NagK, belonging to glycoside hydrolase family 84 (GH84), cleave terminal GlcNAc residues using a substrate-assisted digestion mechanism. AgnC (CpGH89), belonging to glycoside hydrolase family 89 (GH89), both cleaves GlcNAc from the ends of mucin glycans and acts on gastric mucin. Belonging to glycoside hydrolase family 123 (GH123), CpNga123 cleaves GalNAc, but research suggests that it only breaks down glycans taken up by C. perfringens due to the absence of a secretion signal peptide.
Galactosidase
C. perfringens has four galactosidases that belong to the glycoside hydrolase family 2 (GH2): CpGH2A, CpGH2B, CpGH2C, and CpGH2D. Research suggests that these enzymes are effective at breaking down core mucin glycan structures with the ability to bind galactose using CBM51. However, minimal research exists on the specific functioning of galactosidases in C. perfringens.
Fucosidase
Fucose monosaccharides are located on the terminal ends of core O-linked glycans. C. perfringens encodes for three fucosidases that belong to two glycoside hydrolase families: Afc1 and Afc2 in glycoside hydrolase family 29 (GH29), and Afc3 in glycoside hydrolase family 95 (GH95). Afc3 contains a C-terminal CBM51 and is the only fucosidase that contains a carbohydrate-binding module in C. perfringens. Fucosyl residues tend to cover the ends of glycans and protect them against enzymatic digestion, so research suggests that the ability of fucosidases to cleave complex and diverse fucosyl linkages is due to long-term adaptations in C. perfringens that persisted within close range of mucins.
Major toxins
There are five major toxins produced by Clostridium perfringens. Alpha, beta, epsilon and enterotoxin are toxins that increase a cells permeability which causes an ion imbalance while iota toxins destroy the cell's actin cytoskeleton. On the basis of which major, "typing" toxins are produced, C. perfringens can be classified into seven "toxinotypes", A, B, C, D, E, F and G:
Alpha toxin
Alpha toxin (CPA) is a zinc-containing phospholipase C, composed of two structural domains, which destroy a cell's membrane. Alpha toxins are produced by all five types of C. perfringens. This toxin is linked to gas gangrene of humans and animals. Most cases of gas gangrene has been related to a deep wound being contaminated by soil that harbors C. perfringens.
Beta toxin
Beta toxins (CPB) are a protein that causes hemorrhagic necrotizing enteritis and enterotoxaemia in both animals (type B) and humans (type C) which leads to the infected individual's feces becoming bloody and their intestines necrotizing. Proteolytic enzymes, such as trypsin, can break down CPB, making them ineffective. Therefore, the presence of trypsin inhibitors in colostrum makes CPB especially deadly for mammal offspring.
Epsilon toxin
Epsilon toxin (ETX) is a protein produced by type B and type D strains of C. perfringens. This toxin is currently ranked the third most potent bacterial toxin known. ETX causes enterotoxaemia in mainly goats and sheep, but cattle are sometime susceptible to it as well. An experiment using mice found that ETX had an LD50 of 50-110 ng/kg. The excessive production of ETX increases the permeability of the intestines. This causes severe edema in organs such as the brain and kidneys.
The very low LD50 of ETX has led to concern that it may be used as a bioweapon. It appeared on the select agent lists of the US CDC and USDA, until it was removed in 2012. There are no human vaccines for this toxin, but effective vaccines for animals exist.
Iota toxin
Iota toxin (ITX) is a protein produced by type E strains of C. perfringens. Iota toxins are made up of two, unlinked proteins that form a multimeric complex on cells. Iota toxins prevent the formation of filamentous actin. This causes the destruction of the cells cytoskeleton which in turn leads to the death of the cell as it can no longer maintain homeostasis.
Enterotoxin
This toxin (CPE) causes food poisoning. It alters intracellular claudin tight junctions in gut epithelial cells. This pore-forming toxin also can bind to human ileal and colonic epithelium in vitro and necrotize it. Through the caspase-3 pathway, this toxin can cause apoptosis of affected cells. This toxin is linked to type F strains, but has also been found to be produced by certain types of C, D, and E strains.
Other toxins
TpeL is a toxin found in type B, C, and G strains. It is in the same protein family as C. difficile toxin A. It does not appear important in the pathogenesis of types B and C infections, but may contribute to virulence in type G strains. It glycosylates Rho and Ras GTPases, disrupting host cell signaling.
Infection
Tissue necrosis, bacteremia, emphysematous cholecystitis, and gas gangrene, also known as clostridial myonecrosis, have been linked to infections associated with C. perfringens. Research suggests that C. perfringens is capable of engaging in polymicrobial anaerobic infections.
Clostridium perfringens is a common cause of food poisoning in the United States. C. perfringens produces spores, and when these spores are consumed, they produce a toxin that causes diarrhea. Foods cooked in large batches and held at unsafe temperatures (between 40°F and 140°F) are the source of C. perfringens food poisoning outbreaks. Meats such as poultry, beef, and pork are commonly linked to C. perfringens food poisoning. C. perfringens can proliferate in foods that are improperly stored due to the spore's ability to survive normal cooking temperatures. The type A toxin of C. perfringens, also known as the CPA is responsible for food poisoning.
Clostridium perfringens is the most common bacterial agent for gas gangrene. Gas gangrene is induced by α-toxin that embeds itself into the plasma membrane of cells and disrupts normal cellular function by altering membrane structure. Some symptoms include blisters, tachycardia, swelling, and jaundice.
C. perfringens is most commonly known for foodborne illness but can translocate from a gastrointestinal source into the bloodstream which causes bacteremia. C. perfringens bacteremia can lead to toxin-mediated intravascular hemolysis and septic shock. This is rare as it makes up less than 1% of bloodstream isolates but is highly fatal with a reported mortality rate of 27% to 58%.
Clostridium perfringens food poisoning can also lead to another disease known as enteritis necroticans or clostridial necrotizing enteritis, (also known as pigbel); this is caused by C. perfringens type C. This infection is often fatal. Large numbers of C. perfringens grow in the intestines and secrete exotoxin. This exotoxin causes necrosis of the intestines, varying levels of hemorrhaging, and perforation of the intestine. Inflammation usually occurs in sections of the jejunum, midsection of the small intestine. Perfringolysin O (pfoA)-positive C. perfringens strains were also associated with the rapid onset of necrotizing enterocolitis in preterm infants.
A strain of C. perfringens might be implicated in multiple sclerosis (MS) nascent (Pattern III) lesions. Tests in mice found that two strains of intestinal C. perfringens that produced epsilon toxins (ETX) caused MS-like damage in the brain, and earlier work had identified this strain of C. perfringens in a human with MS. MS patients were found to be 10 times more immune-reactive to the epsilon toxin than healthy people.
Tissue gas occurs when C. perfringens infects corpses. It causes extremely accelerated decomposition and can only be stopped by embalming the corpse. Tissue gas most commonly occurs to those who have died from gangrene, large decubitus ulcers, necrotizing fasciitis or to those who had soil, feces, or water contaminated with C. perfringens forced into an open wound.
Clinical Manifestations
Clostridium perfringens infections can lead to various clinical manifestations, ranging from mild gastrointestinal symptoms to life-threatening conditions. The most common presentation is food poisoning, characterized by acute abdominal pain, diarrhea, and, in some cases, vomiting, typically occurring 6 to 24 hours after the ingestion of contaminated food. Unlike many other foodborne illnesses, fever is usually absent. Symptoms are usually self-limiting and resolve within 24 to 48 hours; however, severe dehydration can occur in cases of significant fluid loss. Symptoms of dehydration include dry mouth, decreased urine output, dizziness, and fatigue. Severe symptoms such as diarrhea that persists for more than 48 hours, the inability to keep fluids down, or signs of severe dehydration may necessitate medical attention. Most people are able to recover from C. perfringens food poisoning without treatment. However, people experience diarrhea are usually instructed to drink water or rehydration solutions.
Gas gangrene caused by Clostridium perfringens is characterized by severe symptoms, including intense pain at the injury site, fever, rapid heart rate, sweating, and anxiety. The affected area may show signs of swelling, discoloration (ranging from pale to dark red or purplish), and large, discolored blisters filled with foul-smelling fluid. As the toxins spread, skin and muscle tissue are rapidly destroyed, leading to large areas of dead tissue, gas pockets under the skin (crepitus), and possible renal failure due to red blood cell destruction. Sepsis and septic shock may also occur, which can be fatal.
Necrotizing enteritis caused by Clostridium perfringens presents with a wide range of symptoms, which can vary in severity. The clinical signs range from mild diarrhea to more severe manifestations such as intense abdominal pain, vomiting, bloody stools, and even septic shock. In the most serious cases, the infection can lead to death.
Diagnosis
The diagnosis of Clostridium perfringens food poisoning relies on laboratory detection of the bacterium or its toxin in either a patient’s stool sample or contaminated food linked to the illness. A positive stool culture would have growth of at least 10 cfu/g of C. perfringens. Stool studies include WBCs, ova, and parasites in order to rule out other potential etiologies. ELISA testing is used to detect the CPA toxin. Diagnosing C. perfringens food poisoning is relatively uncommon for several reasons. Most individuals with this foodborne illness do not seek medical care or submit a stool sample for testing, and routine testing for C. perfringens is not typically performed in clinical laboratories. Additionally, public health laboratories generally conduct testing for this pathogen only in the event of an outbreak.
The diagnosis of gas gangrene typically involves several methods to confirm the infection. Imaging techniques such as X-rays, CT scans, or MRIs can reveal gas bubbles or tissue changes indicative of muscle damage. Additionally, bacterial staining or culture of fluid taken from the wound helps identify Clostridium perfringens and other bacteria responsible for the infection. In some cases, a biopsy is performed, where a sample of the affected tissue is analyzed for signs of damage or necrosis.
The diagnosis of clostridial necrotizing enteritis is primarily based on the patient's clinical symptoms, which can include severe abdominal pain, vomiting, and bloody diarrhea. Additionally, confirmation of the presence of Clostridium perfringens type C toxin in stool samples is crucial for accurate diagnosis.
Epidemiology
Clostridium perfringens is responsible for an estimated 966,000 cases annually, or about 10.3% of all foodborne illnesses in which a pathogen is identified. Transmission typically occurs when food contaminated with C. perfringens spores is consumed, allowing the bacteria to produce a toxin in the intestines that causes diarrhea. Outbreaks are often associated with foods cooked in large batches, such as poultry, meat, and gravy, and held at unsafe temperatures between 40-140°F, which allows the bacteria to thrive. These outbreaks tend to occur in settings where large groups are served, such as hospitals, school cafeterias, prisons, nursing homes, and catered events. In most cases, C. perfringens infection causes mild symptoms, including watery diarrhea and mild abdominal cramps, with symptoms typically appearing 8 to 12 hours after consuming contaminated food and resolving within 24 hours. About 90% of affected individuals recover without seeking medical attention, usually within two days. However, vulnerable groups such as the elderly, young children, and immunocompromised individuals face a higher risk of severe complications like dehydration, which can lead to more serious illness or, in rare cases, death. Each year, C. perfringens infections result in approximately 438 hospitalizations and 26 deaths, accounting for 0.8% of foodborne illness-related hospitalizations and 1.9% of associated deaths. Outbreaks are most common in November and December, coinciding with holiday foods like turkey and roast beef. The economic burden of C. perfringens is significant, estimated at $342.7 million annually, including $53.2 million in medical costs, $64.3 million in productivity loss, and $225 million related to fatalities.
Clostridial necrotizing enteritis is rare in the United States; typically, it occurs in populations with a higher risk. Data show that of the 9.4 million cases of foodborne illness in the United States each year, only about 11% are caused by Clostridium perfringens. "Risk factors for enteritis necroticans include protein-deficient diet, unhygienic food preparation, sporadic feasts of meat (after long periods of a protein-deficient diet), diets containing large amounts of trypsin inhibitors (sweet potatoes), and areas prone to infection of the parasite Ascaris (produces a trypsin inhibitor). This disease is contracted in populations living in New Guinea, parts of Africa, Central America, South America, and Asia.
Risk factors for gas gangrene include severe injuries, abdominal surgeries, and underlying health conditions such as colon cancer, diseases of the blood vessels, diabetes, and diverticulitis. However, the most common way to get gas gangrene is through a traumatic injury. In the United States, there is only about 1000 cases of gas gangrene per year. When addressed with adequate care, gas gangrene has a mortality rate of 20-30% but has a mortality rate of 100% if left untreated.
Food poisoning incidents
On May 7, 2010, 42 residents and 12 staff members at a Louisiana (USA) state psychiatric hospital were affected and experienced vomiting, abdominal cramps, and diarrhea. Three patients died within 24 hours. The outbreak was linked to chicken which was cooked a day before it was served and was not cooled down according to hospital guidelines. The outbreak affected 31% of the residents of the hospital and 69% of the staff who ate the chicken. How many of the affected residents ate the chicken is unknown.
In May 2011, a man died after allegedly eating food contaminated with the bacteria on a transatlantic American Airlines flight. The man's wife and daughter were suing American and LSG Sky Chefs, the German company that prepared the inflight food.
In December 2012, a 46-year-old woman died two days after eating a Christmas Day meal at a pub in Hornchurch, Essex, England. She was among about 30 people to fall ill after eating the meal. Samples taken from the victims contained C. perfringens. The hotel manager and the cook were jailed for forging cooking records relating to the cooking of the turkey.
In December 2014, 87-year-old Bessie Scott died three days after eating a church potluck supper in Nackawic, New Brunswick, Canada. Over 30 other people reported signs of gastrointestinal illness, diarrhea, and abdominal pain. The province's acting chief medical officer says, Clostridium perfringens is the bacteria [sic] that most likely caused the woman's death.
In October 2016, 66-year-old Alex Zdravich died four days after eating an enchilada, burrito, and taco at Agave Azul in West Lafayette, Indiana, United States. Three others who dined the same day reported signs of foodborne illness, which were consistent with the symptoms and rapid onset of C. perfringens infection. They later tested positive for the presence of the bacteria, but the leftover food brought home by Zdravich tested negative.
In November 2016, food contaminated with C. perfringens caused three individuals to die, and another 22 to be sickened, after a Thanksgiving luncheon hosted by a church in Antioch, California, United States.
In January 2017, a mother and her son sued a restaurant in Rochester, New York, United States, as they and 260 other people were sickened after eating foods contaminated with C. perfringens. "Officials from the Monroe County Department of Public Health closed down the Golden Ponds after more than a fourth of its Thanksgiving Day guests became ill. An inspection revealed a walk-in refrigerator with food spills and mold, a damaged gasket preventing the door from closing, and mildew growing inside."
In July 2018, 647 people reported symptoms after eating at a Chipotle Mexican Grill restaurant in Powell, Ohio, United States. Stool samples tested by the CDC tested positive for C. perfringens.
In November 2018, approximately 300 people in Concord, North Carolina, United States, were sickened by food at a church barbecue that tested positive for C. perfringens.
In 2021, dozens of hospital workers in Alaska were sick and it was traced back to a Cubano Sandwich. Health officials wrote that almost all symptoms resolved within 24 hours. No one who ate the food reportedly needed hospitalization. It is rare for Alaska to see an outbreak with this magnitude when it's not associated with some sort of national food borne illness.
Prevention
Preventing Clostridium perfringens contamination and growth involves careful food handling, proper cooking, and appropriate storage practices. Most foods, especially beef and chicken, can be protected by cooking them to the recommended internal temperatures. Using a kitchen thermometer is the most reliable way to check that meats reach safe cooking temperatures. As a general rule, food should be avoided if it smells, tastes, looks off, or has been left out at unsafe temperatures for a long period of time.
C. perfringens spores can multiply within a temperature range of 59°F (15°C) to 122°F (50°C). To prevent bacterial growth, leftovers should be refrigerated within two hours of preparation, with their temperature chilled down to below 40°F (4°C). Large portions of food that contain meat, should be divided into smaller containers before refrigeration to ensure even cooling. Before serving leftovers, they should be reheated to at least 165°F (74°C) to destroy any bacteria that may have grown during storage.
High-risk foods, such as canned vegetables, smoked or cured meats, and salted or smoked fish, require additional attention. Improper processing or storage can allow bacteria to grow and produce dangerous toxins. Signs of contamination, such as unusual odors, changes in texture, or bulging cans (also known as "bombage"), indicate food spoilage and should be disposed.
Preventing gas gangrene involves taking precautions to avoid bacterial infections. Healthcare providers follow strict protocols to prevent infections, including those caused by Clostridium perfringens. To reduce the risk of gas gangrene, individuals should clean wounds thoroughly with soap and water and seek medical attention for deep or uncleanable wounds. It is also essential to monitor injuries for changes in skin condition or the onset of severe pain. Wearing protective gear when engaging in activities like biking or motorcycling can help prevent injury. Additionally, working with healthcare providers to manage underlying conditions that affect circulation or weaken the immune system can further reduce the risk of infection.
Treatment
The treatment of Clostridium perfringens infections depends on the type and severity of the condition. For severe infections, such as gas gangrene (clostridial myonecrosis), the primary approach involves surgical debridement of the affected area. This procedure removes devitalized tissue where bacteria grow, which limits the spread of the infection. Antimicrobial therapy is usually started at the same time, with penicillin being the most commonly used drug. However, C. perfringens shows different resistance patterns with about 20% of strains being resistant to clindamycin, and 10% being resistant to metronidazole. C. perfringens is often more susceptible to vancomycin when compared to other pathogenic Clostridia, making it an alternative option for treatment in some cases.
Therapies, such as hyperbaric oxygen therapy (HBOT), may also be used for severe clostridial tissue infections. HBOT increases oxygen delivery to infected tissues, creating an environment that inhibits the growth of anaerobic bacteria like C. perfringens. While not commonly used, HBOT can be beneficial in certain cases.
For foodborne illness caused by C. perfringens, treatment is typically unnecessary. Most people who suffer from food poisoning caused by C. perfringens usually fight off the illness without the need of any antibiotics. Extra fluids should be drank consistently until diarrhea dissipates.
Research
C. perfringens has shown increasing multidrug resistance, particularly in strains from humans and animals. High resistance levels were found with antibiotics such as tetracycline, erythromycin, and sulfonamides. Genetic factors, misuse of antibiotics, and bacterial evolution are the cause of this issue. This highlights the importance of finding new treatment strategies.
Multilocus Sequence Typing (MLST) and Whole Genome Sequencing (WGS) have been used to find the genetic diversity of C. perfringens. These methods have identified 195 distinct sequence types grouped into 25 clonal complexes from 322 genomes. Phylogenetic groups were also found in multiple different hosts and environmental sources. This highlights the bacteria's transmission potential and adaptability across species.
| Biology and health sciences | Gram-positive bacteria | Plants |
475779 | https://en.wikipedia.org/wiki/Calliphoridae | Calliphoridae | The Calliphoridae (commonly known as blowflies, blow flies, blow-flies, carrion flies, bluebottles, or greenbottles) are a family of insects in the order Diptera, with almost 1,900 known species. The maggot larvae, often used as fishing bait, are known as gentles. The family is known to be polyphyletic, but much remains disputed regarding proper treatment of the constituent taxa, some of which are occasionally accorded family status (e.g., Bengaliidae and Helicoboscidae).
Description
Characteristics
Calliphoridae adults are commonly shiny with metallic colouring, often with blue, green, or black thoraces and abdomens. Antennae are three-segmented and aristate. The aristae are plumose their entire length, and the second antennal segment is distinctly grooved. Members of Calliphoridae have branched Rs 2 veins, frontal sutures are present, and calypters are well developed.
The characteristics and arrangements of hairlike bristles are used to differentiate among members of this family. All blowflies have bristles located on the meron. Having two notopleural bristles and a hindmost posthumeral bristle located lateral to presutural bristle are characteristics to look for when identifying this family.
The thorax has the continuous dorsal suture across the middle, along with well-defined posterior calli. The postscutellum is absent or weakly developed. The costa is unbroken and the subcosta is apparent on the insect.
Development
Most species of blowflies studied thus far are anautogenous; a female requires a substantial amount of protein to develop mature eggs within her ovaries (about 800 μg per pair of ovaries in Phormia regina). The current theory is that females visit carrion both for protein and egg laying, but this remains to be proven. Blowfly eggs, usually yellowish or white in color, are about 1.5 mm × 0.4 mm, and when laid, look like rice grains. While the female blowfly typically lays 150–200 eggs per batch, she is usually iteroparous, laying around 2,000 eggs during the course of her life. The sex ratio of blowfly eggs is usually 50:50, but one exception is females from two species of the genus Chrysomya (C. rufifacies and C. albiceps), which are either arrhenogenic (laying only male offspring) or thelygenic (laying only female offspring).
Hatching from an egg to the first larval stage takes about 8 hours to a day. Larvae have three stages of development (instars); each stage is separated by a molting event. The instars are separable by examining the posterior spiracles, or openings to the breathing system. The larvae use proteolytic enzymes in their excreta (as well as mechanical grinding by mouth hooks) to break down proteins on the livestock or corpse on which they are feeding. Blowflies are poikilothermic – the rate at which they grow and develop is highly dependent on temperature and species. Under room temperature (about 20 °C), the black blowfly Phormia regina can change from egg to pupa in 150–266 hours (six to 11 days). When the third larval stage is complete, it leaves the corpse and burrows into the ground to pupate, emerging as an adult 7–14 days later.
Food sources
Adult blowflies are occasional pollinators, being attracted to flowers with strong odors resembling rotting meat, such as the American pawpaw or dead horse arum. Little doubt remains that these flies use nectar as a source of carbohydrates to fuel flight, but just how and when this happens is unknown. One study showed the visual stimulus a blowfly receives from its compound eyes is responsible for causing its legs to extend from its flight position and allow it to land on any surface.
Larvae of most species are scavengers of carrion and dung, and most likely constitute the majority of the maggots found in such material, although they are not uncommonly found in close association with other dipterous larvae from the families Sarcophagidae and Muscidae, and many other acalyptrate muscoid flies.
Predators
Predators of blowflies include spiders, beetles, frogs, and birds, including chickens.
In the Chihuahuan desert of Mexico, a fungus, Furia vomitoriae (from the family of Entomophthoraceae) affects bluebottle flies. It forms masses of conidiophores erupting through the intersegmental areas (or clear bands) on the abdominal dorsum of the flies and eventually kills them.
Diversity
About 1,900 species of blowflies are known, with 120 species in the Neotropics, and a large number of species in Africa and Southern Europe.
Their typical habitats are temperate to tropical areas that provide a layer of loose, damp soil and litter where larvae may thrive and pupate.
Genera
Sources: MYIA, FE, Nomina, A/O DC
This is a selected list of genera from the Palearctic, Nearctic, Malaysia (Japan), and Australasia:
Abago Grunin, 1966
Amenia Robineau-Desvoidy, 1830
Angioneura Brauer & Bergenstamm, 1893
Apaulina Hall, 1948
Cynomya Robineau-Desvoidy, 1830
Aphyssura Hardy, 1940
Auchmeromyia Brauer & Bergenstamm, 1891
Bellardia Robineau-Desvoidy, 1863
Bengalia Robineau-Desvoidy, 1830
Booponus Aldrich, 1923
Boreellus Aldrich & Shannon, 1923
Caiusa Surcouf, 1920
Calliphora Robineau-Desvoidy, 1830
Callitroga Hall, 1948
Catapicephala Macquart, 1851
Chloroprocta Wulp, 1896
Chrysomya Robineau-Desvoidy, 1830
Cochliomyia Townsend, 1915
Compsomyiops Townsend, 1918
Cordylobia Gruenberg, 1903
Cyanus Hall, 1948
Dyscritomyia Grimshaw, 1901
Eggisops Rondani, 1862
Eucalliphora Townsend, 1908
Eumesembrinella Townsend, 1931
Eurychaeta Brauer & Bergenstamm, 1891
Euphumosia Malloch, 1926
Hemilucilia Brauer, 1895
Hemipyrellia Townsend, 1918
Lucilia Robineau-Desvoidy, 1830
Melanomya Rondani, 1856
Melinda Robineau-Desvoidy, 1830
Mufetiella Villeneuve, 1933
Nesodexia Villeneuve, 1911
Neta Shannon, 1926
Onesia Robineau-Desvoidy, 1830
Opsodexia Townsend, 1915
Pachychoeromyia Villeneuve, 1920
Paralucilia Brauer & Bergenstamm, 1891
Paramenia Brauer & Bergenstamm, 1889
Paraplatytropesa Crosskey, 1965
Phormia Robineau-Desvoidy, 1830
Phumosia Robineau-Desvoidy, 1830
Platytropesa Macquart, 1851
Polleniopsis Townsend, 1917
Prosthetosoma Silvestri, 1920
Protocalliphora Hough, 1899
Protophormia Townsend, 1908
Ptilonesia Bezzi, 1927
Rhynchoestrus Séguy, 1926
Sarconesia Bigot, 1857
Silbomyia Macquart, 1843
Stilbomyella Malloch, 1935
Toxotarsus Macquart, 1851
Triceratopyga Rohdendorf, 1931
Tricyclea Wulp, 1885
TricycleopsisVilleneuve, 1927
Trypocalliphora Peus, 1960
Xenocalliphora Malloch, 1924
Economic importance
Myiasis
Blowflies have caught the interest of researchers in a variety of fields, although the large body of literature on calliphorids has been concentrated on solving the problem of myiasis in livestock. The sheep blowfly Lucilia cuprina causes the Australian sheep industry an estimated AU$170 million a year in losses.
The most common causes of myiasis in humans and animals are the three dipteran families Oestridae, Calliphoridae, and Sarcophagidae. Myiasis in humans is clinically categorized in six ways: dermal and subdermal, facial cavity, wound or trauma, gastrointestinal, vaginal, and generalized. If found in humans, the dipteran larvae are usually in their first instar. The only treatment necessary is just to remove the maggots, and the patient heals naturally. Whilst not strictly a myiasis species, the Congo floor maggot feeds on mammal blood, occasionally human.
Screwworms
The New World primary screwworm (Cochliomyia hominivorax), once a major pest in Southern United States, has been eradicated from the United States, Mexico, and Central America through an extensive release program by the USDA of sterilized males. The USDA maintains a sterile screwworm fly production plant and release program in the eastern half of the Republic of Panama to keep fertile screwworms from migrating north. Currently, this species is limited to lowland tropical countries in South America and some Caribbean islands.
The Old World primary screwworm (Chrysomya bezziana) is an obligate parasite of mammals. This fly is distributed throughout the Old World, including Southeast Asia, tropical and subtropical Africa, some countries in the Middle East, India, the Malay Peninsula, the Indonesian and Philippine Islands, and Papua New Guinea.
The secondary screwworm (Cochliomyia macellaria) has become one of the principal species on which to base post mortem interval estimations because its succession and occurrence on decomposing remains has been well defined. The secondary screwworm is found throughout the United States and the American tropics, and in southern Canada during summers. This species is one of the most common species found on decomposing remains in the US South.
Maggot therapy
Maggot debridement therapy (MDT) is the medical use of selected, laboratory-raised fly larvae for cleaning nonhealing wounds. Medicinal maggots perform debridement by selectively eating only dead tissue. Lucilia sericata (Phaenicia sericata), or the common green bottlefly, is the preferred species used in maggot therapy. MDT can be used to treat pressure ulcers, diabetic foot wounds, venous stasis ulcers, and postsurgical wounds.
Disease
Adults may be vectors of pathogens of diseases such as dysentery. Flies, most commonly Calliphoridae, have frequently been associated with disease transmission in humans and animals, as well as myiasis. Studies and research have linked Calliphora and Lucilia to vectors of causal agents of bacterial infections. These larvae, commonly seen on decaying bodies, feed on carrion while the adults can be necrophagous or vegetative. During the process of decay, microorganisms (e.g. Mycobacterium) may be released through the body. Flies arrive at the scene and lay their eggs. The larvae begin eating and breaking down the corpse, simultaneously ingesting these organisms which is the first step of one transmission route.
The bacterium which causes paratuberculosis in cattle, pigs and birds (M. a. avium) has been isolated and recovered from these flies through several different experiments.
Other potential and threatening diseases include rabbit haemorrhagic disease in New Zealand and flystrike. Although strike is not limited to blow flies, these maggots are a major source of this skin invasion, causing lesions, which, if severe enough, may be lethal. Strike starts when blow flies lay eggs in a wound or fecal material present on the sheep. When the maggots hatch, they begin feeding on the sheep and thus irritating it. As soon as the first wave of maggots hatch, they attract more blow flies, causing the strike. Insecticides are available for blow fly prevention (typically containing cypermethrin), and precautionary measures may be taken, such as docking tails, shearing, and keeping the sheep healthy overall.
Salmonellosis has also been proven to be transmitted by the blow fly through saliva, feces and direct contact by the flies' tarsi. Adult flies may be able to spread pathogens via their sponging mouthparts, vomit, intestinal tract, sticky pads of their feet, or even their body or leg hairs.
As the flies are vectors of many diseases, the importance of identifying the transmissible agents, the route of transmission, and prevention and treatments in the event of contact are becoming increasingly important. With the ability to lay hundreds of eggs in a lifetime and the presence of thousands of larvae at a time in such close proximity, the potential for transmission is high, especially at ideal temperatures.
Pollination
Calliphoridae are, alongside managed and wild bees, likely to be the main crop pollinating insect. They visit (and thus may pollinate) flowers of a wide range of plants, including crop plants (e.g. avocado, mango, onion, leek, carrot, cauliflower). Their sponging mouthparts mean that when visiting flowers, their head and upper body must broadly contact the inside of the flower. They have numerous hairs, including on the head and thorax, which may help them carry pollen, and indeed calliphorids in the wild have been observed carrying large amounts of pollen. Compared to honey bees, blow flies are active under a broader range of environmental conditions. However, it is unknown how their pollination abilities compare to those of bees, there are few studies assessing their contribution to pollination, and the exact species that pollinate are often not identified.
Forensic importance
Blow flies are usually the first insects to come in contact with carrion because they have the ability to smell dead animal matter from up to away. Upon reaching the carrion, females deposit eggs on it. Since development is highly predictable if the ambient temperature is known, blow flies are considered a valuable tool in forensic science. Blow flies are used forensically to estimate the minimum post mortem interval (PMImin) for human corpses. Traditional estimations of time since death are generally unreliable after 72 hours and often entomologists are the only officials capable of generating an accurate approximate time interval. The specialized discipline related to this practice is known as forensic entomology.
In addition to being used to estimate the PMImin, assuming colonization occurred after death, blow fly specimens found infesting a human corpse are used to determine if the corpse was relocated or if the individual ingested narcotics prior to death.
Calliphora vicina and Cynomya mortuorum are important flies of forensic entomology. Other forensically important Calliphoridae are Phormia regina, Calliphora vomitoria, Calliphora livida, Lucilia cuprina, Lucilia sericata, Lucilia illustris, Chrysomya rufifacies, Chrysomya megacephala, Cochliomyia macellaria, and Protophormia terraenovae. One myth states that species from the genus Lucilia can sense death and show up right before it even occurs.
| Biology and health sciences | Flies (Diptera) | null |
475861 | https://en.wikipedia.org/wiki/Jalape%C3%B1o | Jalapeño | The jalapeño ( , , ) is a medium-sized chili pepper pod type cultivar of the species Capsicum annuum. A mature jalapeño chili is long and wide, and hangs down from the plant. The pungency of jalapeño peppers varies, but is usually between 4,000 and 8,500 units on the Scoville scale. Commonly picked and consumed while still green, it is occasionally allowed to fully ripen and turn red, orange, or yellow. It is wider and generally milder than the similar Serrano pepper.
History and etymology
The jalapeño is variously named huachinango, for the ripe red jalapeño, and chile gordo (meaning "fat chili pepper") also known as cuaresmeño.
The name jalapeño is Spanish for "from Xalapa", the capital city of Veracruz, Mexico, where the pepper was traditionally cultivated.
Genetic analysis of Capsicum annuum places jalapeños as a distinct genetic clade with no close sisters that are not directly derived from jalapeños. Jalapeños were in use by the Aztecs prior to the Spanish conquest; Bernardino de Sahagún in the Florentine Codex writes of Aztec markets selling chipotles (smoked jalapeños) and mole made from chipotles, besides the sale of fresh chilies. The use of peppers in the Americas dates back thousands of years, including the practice of smoking some varieties of peppers in order to preserve them; further well preserved samples and genetic testing would be needed to determine the usage and existence of the jalapeño clade and pod type into the past.
Cultivation
In 1999, roughly of land in Mexico was dedicated to jalapeño production; , that had decreased to . Jalapeños account for thirty percent of Mexico's chili production, and while the total land area used for cultivation has decreased, there has been a 1.5% increase in volume yield per year in Mexico due to increasing irrigation, use of greenhouses, better equipment, knowledge, and improved techniques. Because of this, in 2009, 619,000 tons of jalapeños were produced with 42% of the crop coming from Chihuahua, 12.9% from Sinaloa, 6.6% from Jalisco, and 6.3% from Michoacán. La Costeña controls about 60% of the world market and, according to company published figures, exports 16% of the peppers that Mexico produces, an 80% share of the 20% that Mexico exports in total. The US imports 98% of La Costeña's exports.
According to the USDA, since 2010 California produces the most jalapeños, followed by New Mexico and Texas - a total of of peppers in 2014. It is difficult to get accurate statistics on chilies and specific chilies as growers are not fond of keeping and sharing such data and reporting agencies often lump all green chilies together, or all hot chilies, with no separation of pod type. In New Mexico in 2002 the crop of jalapeños were worth $3 million at the farm gate and $20 million after processing.
China, Peru, Spain, and India also produce commercial chilies, including jalapeños.
Jalapeños are a pod type of Capsicum annuum. The growing period is 70–80 days. When mature, the plant stands tall. Typically, a plant produces 25 to 35 pods. During a growing period, a plant will be picked multiple times. As the growing season ends, the peppers turn red, as seen in sriracha sauce. Jalapeños thrive in a number of soil types and temperatures, though they prefer warmer climates, provided they have adequate water. The optimum temperature for seed germination is , with degradation of germination seen above and little to no germination occurring at ; at the time to 50% germination rate depends on cultivar and seed lot but was tested as being between 4 and 5 days, which is shorter than cayenne. A pH of 4.5 to 7.0 is preferred for growing jalapeños, and well-drained soil is essential for healthy plants. Jalapeños need at least 6 to 8 hours of sunlight per day. Experiments show that unlike bell peppers at least 7.5 millimolar (mM) nitrogen is needed for optimal pod production, and 15 to 22 mM nitrogen produces the best result: the plant produces both more leaves and more pods, rather than just more leaves. Once picked, individual peppers may turn to red of their own accord. The peppers can be eaten green or red. Though usually grown as an annual they are perennial and if protected from frost can produce during multiple years, as with all Capsicum annuum.
Jalapeños are subject to root rot and foliar blight, both often caused by Phytophthora capsici; over-watering worsens the condition as the fungus grows best in warm wet environments. Crop rotation can help, and resistant strains of jalapeño, such as the 'NuMex Vaquero' and 'TAM Mild Jalapeño', have been and are being bred as this is of major commercial impact throughout the world. As jalapeños are a cultivar, the diseases are common to Capsicum annuum: Verticillium wilt, Cercospora capsici, Powdery mildew, Colletotrichum capsici (Ripe Rot), Erwinia carotovora (Soft Rot), Beet curly top virus, Tospovirus (Tomato spotted wilt virus), Pepper mottle virus, Tobacco mosaic virus, Pepper Geminiviridae, and Root-knot nematode being among the major commercially important diseases.
After harvest, if jalapeños are stored at they have a shelf life of up to 3–5 weeks. Jalapeños produce 0.1–0.2 μL per kg per hour of ethylene, very low for chilies, and do not respond to ethylene treatment. Holding jalapeños at 20–25 °C and high humidity can be used to complete the ripening of picked jalapeños. A hot water dip of for 4 minutes is used to kill off molds that may exist on the picked peppers without damaging them. The majority of jalapeños are wet processed, canned, or pickled on harvesting for use in mixes, prepared food products, and salsas.
Hybrids and sub-cultivars
There are a wide variety of breeds for consumer and commercial use of jalapeño plants. The majority fall under one of four categories: F1 hybrids, where the parent plants have been hand-emasculated and cross-bred to produce uniform offspring with hybrid vigor; cultivars which are F-11 or F-12 hybrids or later generations where a stable unique population has been developed; landraces; and F2 hybrids.
F1 hybrids produce the highest and most uniform yields but cost 25 times the cost of open-pollinated seed, leading to only 2% of the farmland dedicated to jalapeño cultivation in the United States being planted with F1 hybrids. F2 hybrids often produce similarly to F1 hybrids; however, some F1 hybrids are produced via recessive male sterility to eliminate the need to hand-pollinate, reducing the cost to produce the hybrid, but producing a 25% reduction in yield in the F2 generation. Some notable F1 hybrids are 'Mitla', 'Perfecto', 'Tula', 'Grande' (a hot jalapeño), 'Sayula', 'Senorita', and 'Torreon', most of them being developed and marketed by Petoseed, a brand of Seminis.
Cultivars are researched and created to promote desirable traits. Common traits selected for are resistance to viruses and other pepper-related diseases, milder peppers, early ripening, more attractive fruit in terms of size, wall thickness, and corking, and higher yields. The land-grant universities and the Chile Pepper Institute promote the use of cultivars as the most sustainable and environmentally safe disease control method both in terms of economics and long-term environmental perspective. Notable cultivars include 'Early Jalapeño', 'TAM Mild Jalapeño', 'TAM Mild Jalapeño II', 'TAM Veracruz', the yellow 'TAM Jaloro', 'NuMex Vaquero', the colorful 'NuMex Piñata', 'TAM Dulcito', 'Waialua', Biker Billy and 'NuMex Primavera'.
Sweet hybrids
Sweet hybridized varieties have been created with no "heat", although they retain the look and flavor of a jalapeño.
Eating characteristics
Nutrients
A raw jalapeño is 92% water, 6% carbohydrates, 1% protein, and contains negligible fat (table). A reference serving of raw jalapeños provides of food energy, and is a rich source (20% or more of the Daily Value, DV) of vitamin C, vitamin B6, and vitamin E, with vitamin K in a moderate amount (table). Other micronutrients are low in content (table).
Scoville heat units
Compared with other chillies, the jalapeño heat level varies from mild to hot depending on cultivation and preparation and can have from a few thousand to over 10,000 Scoville heat units. The number of scars on the pepper, which appear as small brown lines, called 'corking', has a positive correlation with heat level, as growing conditions which increase heat level also cause the pepper to form scars. For US consumer markets, 'corking' is considered unattractive; however, in other markets, it is a favored trait, particularly in pickled or oil-preserved jalapeños.
The heat level of jalapeños varies even for fruit from the same plant; however, some cultivars have been bred to be generally milder, and on the low side of the heat range, such as the 'TAM Milds' and 'Dulcito', and others to be generally hotter, and on the high end of the heat range, such as 'Grande'. As the peppers ripen their pungency increases, making red jalapeños to be generally hotter than green jalapeños, at least of the same variety. If the jalapeño plants were stressed by increased water salinity, erratic watering, temperature, light, soil nutrition, insects, or illness, this will increase their pungency.
All of the capsaicin and related compounds are concentrated in vesicles found in the placenta membrane surrounding the seeds; the vesicles appear white or yellow and fluoresce in the range of 530 to 600 nm when placed in violet light. If fresh chili peppers come in contact with the skin, eyes, lips or other membranes, irritation can occur; some people who are particularly sensitive wear latex or vinyl gloves while handling peppers. If irritation does occur, washing the oils off with hot soapy water and applying vegetable oil to the skin may help. When preparing jalapeños, it is recommended that hands not come in contact with the eyes as this leads to burning and redness.
Serving methods
Stuffed jalapeños are hollowed-out fresh jalapeños (served cooked or raw) filled with seafood, meat, poultry, or cheese.
Pickled jalapeños, a type of pickled pepper, sliced or whole, are often served hot or cold on top of nachos, which are tortilla chips with melted cheese on top, a Tex-Mex dish.
Chipotles are smoked ripe jalapeños.
Jalapeño jelly, which is a pepper jelly, can be prepared using jelling methods.
Jalapeño peppers are often muddled and served in mixed drinks.
Jalapeño poppers are an appetizer; jalapeños are stuffed with cheese, usually cheddar or cream cheese, breaded or wrapped in bacon, and cooked.
Armadillo eggs are jalapeños or similar chilis stuffed with cheese, coated in seasoned sausage meat and wrapped in bacon. The "eggs" are then grilled until the bacon starts to crisp.
Chiles toreados are fresh jalapeños that are sauteed in oil until the skin is blistered all over. They are sometimes served with melted cheese on top.
Texas toothpicks are jalapeños and onions shaved into straws, lightly breaded, and deep-fried.
Chopped jalapeños are a common ingredient in many salsas and chilis.
Jalapeño slices are commonly served in Vietnamese pho and bánh mì, and are also a common sandwich and pizza topping in the West.
Culinary concerns
Jalapeños are a low-acid food with a pH of 4.8–6.0 depending on maturity and individual pepper. If canned or pickled jalapeños appear gassy, mushy, moldy, or have a disagreeable odor, then to avoid botulism, special precautions are needed to avoid illness and spread of the bacteria. Canning or packaging in calcium chloride increases the firmness of the peppers and the calcium content, whether or not the peppers are pickled as well as canned.
In 2008, fresh jalapeños from Mexico were tested positive for Salmonella leading the FDA to believe that the peppers were responsible for much of the 2008 United States salmonellosis outbreak. This large outbreak of Salmonella led to increased research into the detection of pathogens on jalapeños, the frequency and behavior of foodborne illness related to jalapeños, and ways to prevent foodborne illnesses from fresh jalapeños. Contaminated irrigation water and processing water are the two most common methods by which jalapeños become contaminated, as was the case in the 2008 outbreak. Jalapeños have similar microbial properties to tomatoes. The outer layer of their skin provides a safe environment for pathogens to survive, and if damaged or chopped provides a growth medium for these pathogens. Washing fresh jalapeños is important to reduce pathogen counts both at the farm and consumer level, but without cold storage it is insufficient to prevent pathogen spread.
In culture
The jalapeño is a Mexican chili but was designated by the Texas Legislature as the official "State Pepper of Texas" in 1995. In Mexico, jalapeños are used in many forms such as in salsa, pico de gallo, or grilled jalapeños. Jalapeños were included as food on the Space Shuttle as early as 1982.
Guinness World Records recognizes Alfredo Hernandes for the most jalapeños eaten in a minute: 16, on 17 September 2006 at the La Costeña Feel the Heat Challenge in Chicago, Illinois. Patrick Bertoletti holds the Major League Eating jalapeño records at 275 pickled jalapeños in 8 minutes on 1 May 2011, and 191 pickled jalapeños in 6.5 minutes on 16 September 2007 in the 'Short-Form'.
Joaquín Guzmán, also known as "El Chapo", the leader of the Sinaloa Cartel, operated a cannery in Guadalajara producing "Comadre Jalapeños" in order to ship cocaine to the US.
National Jalapeño Day is September 5th in the United States.
Gallery
| Biology and health sciences | Botanical fruits used as culinary vegetables | Plants |
476052 | https://en.wikipedia.org/wiki/Beloniformes | Beloniformes | Beloniformes is an order composed of six families (and about 264 species) of freshwater and marine ray-finned fish:
Adrianichthyidae (ricefish and medakas)
Belonidae (needlefish)
Exocoetidae (flyingfishes)
Hemiramphidae (halfbeaks)
Scomberesocidae (sauries)
Zenarchopteridae (viviparous halfbeaks)
With the exception of the Adrianichthyidae, these are streamlined, medium-sized fishes that live close to the surface of the water, feeding on algae, plankton, or smaller animals including other fishes. Most are marine, though a few needlefish and halfbeaks inhabit brackish and fresh waters.
The order is sometimes divided up into two suborders, the Adrianichthyoidei and the Belonoidei, although this clade is referred to as Exocoetoidei in the 5th edition of Fishes of the World. The Adrianichthyoidei contain only a single family, the Adrianichthyidae. Originally, the Adrianichthyidae were included in the Cyprinodontiformes and assumed to be closely related to the killifish, but a closer relationship to the beloniforms is indicated by various characteristics including the absence of the interhyal, resulting in the upper jaw being fixed or not protrusible. The Belonoidei may also be further subdivided into two superfamilies, the Scomberesocoidea and the Exocoetoidea. The Scomberesocoidea contain the Belonidae and Scomberesocidae, while the Exocoetoidea comprise the Exocoetidae, Hemiramphidae and Zenarchopteridae. However, newer evidence shows the flyingfishes are nested within the halfbeaks, and the needlefish and sauries are nested within the subfamily Zenarchopterinae of the family Hemiramphidae, which has been recognized as its own family. The sauries are also nested within the family Belonidae.
The beloniforms display an interesting array of jaw morphologies. The basal condition in the order excluding the ricefishes is an elongated lower jaw in juveniles and adults as represented in halfbeaks. In the needlefish and sauries, both jaws are elongated in the adults; the juveniles of most species develop through a "halfbeak stage" before having both jaws elongated. The elongated lower jaw is lost in adults and is lost in most juveniles in the flyingfishes and some halfbeak genera.
They are known for many commercial uses, and have about 260 different species. Beloniformes lack a complete sequence of mitogenomes. This leads to many variations in mtDNA, about 35 different ones. To understand evolution for Beloniformes and to identify the larvae, scientists will use Beloniformes to help them study this.
Timeline of genera
| Biology and health sciences | Acanthomorpha | Animals |
477011 | https://en.wikipedia.org/wiki/Allium%20ursinum | Allium ursinum | Allium ursinum, known as wild garlic, ramsons, cowleekes, cows's leek, cowleek, buckrams, broad-leaved garlic, wood garlic, bear leek, Eurasian wild garlic or bear's garlic, is a bulbous perennial flowering plant in the amaryllis family, Amaryllidaceae. It is native to Eurasia, where it grows in moist woodland. It is a wild relative of onion and garlic, all belonging to the same genus, Allium. There are two recognized subspecies: A. ursinum subsp. ursinum and A. ursinum subsp. ucrainicum.
Etymology
The Latin specific name ursinum translates to 'bear' and refers to the supposed fondness of the brown bear for the bulbs; folk tales describe the bears consuming them after awakening from hibernation. Another theory is that the "ursinum" may refer to Ursa Major, as A. ursinum was perhaps one of the most northerly distributed Allium species known to the ancient Greeks, though this hypothesis is disputed. Common names for the plant in many languages also make reference to bears.
Cows love to eat them, hence the modern vernacular name of cows's leek. In Devon, dairy farmers have occasionally had the milk of their herds rejected because of the garlic flavour imparted to it by the cows having grazed upon the plant.
Ramsons is from the Old English word hramsa, meaning "garlic". There is evidence it has been used in British cuisine since the Celtic Britons over 1,500 years ago.
Early healers among the Celts, Gaels, and Teutonic tribes and ancient Romans were familiar with the wild herb who called it herba salutaris, meaning 'healing herb'.
Description
Allium ursinum is a bulbous, perennial herbaceous monocot, that reproduces primarily by seed. The narrow bulbs are formed from a single leaf base and produce bright green entire, elliptical leaves up to long x wide with a petiole up to long. The inflorescence is an umbel of six to 20 white flowers, lacking the bulbils produced by some other Allium species such as Allium vineale (crow garlic) and Allium oleraceum (field garlic). The flowers are star-like with six white tepals, about in diameter, with stamens shorter than the perianth.
It flowers in the British Isles from April to June, starting before deciduous trees leaf in the spring. The flower stem is triangular in cross-section and the leaves are broadly lanceolate, similar to those of the toxic lily of the valley (Convallaria majalis).
Similarity to poisonous plants
Plants that may be mistaken for A. ursinum include lily of the valley, Colchicum autumnale, Arum maculatum, and Veratrum viride or Veratrum album, all of which are poisonous. In Europe, where ramsons are popularly harvested from the wild, people are regularly poisoned after mistakenly picking lily of the valley or Colchicum autumnale.
Grinding the leaves between the fingers and checking for a garlic-like smell can be helpful, but if the smell remains on the hands, one can mistake a subsequent poisonous plant for a safe one. When the leaves of A. ursinum and Arum maculatum first sprout, they look similar, but unfolded Arum maculatum leaves have irregular edges and many deep veins, while ramsons leaves are convex with a single main vein. The leaves of lily of the valley are in pairs, dull green, and come from a single reddish-purple stem, while the leaves of A. ursinum each have their own stem, are shiny when new, and are bright green.
Distribution and habitat
It is native to Europe and Asia, where it grows in moist woodland. It can be found in temperate Europe from Ireland east to the Caucasus. It is common in much of the lowlands of the British Isles with the exception of the far north of Scotland, Orkney and Shetland. The ursinum subspecies is found in western and central Europe, while the ucrainicum subspecies is found in the east and southeast.
It grows in deciduous woodlands with moist soils, preferring slightly acidic conditions. In the British Isles, colonies are frequently associated with bluebells (Hyacinthoides non-scripta), especially in ancient woodland. It is considered to be an ancient woodland indicator species.
Ecology
As its name suggests, A. ursinum is an important food for brown bears. The plant is also a favourite of wild boar.
A. ursinum is the primary larval host plant for a specialised hoverfly, Portevinia maculata (ramsons hoverfly).
The flowers are pollinated by bees.
Uses
All parts of A. ursinum are edible. The leaves can be used as salad, herb, boiled as a vegetable, in soup, or as an ingredient for a sauce that may be a substitute for pesto in lieu of basil. Leaves are also often used to make garlic butter. In Russia the stems are preserved by salting and eaten as a salad. A variety of Cornish Yarg cheese has a rind coated in wild garlic leaves. The leaves can be pickled in the same way as Allium ochotense known as mountain garlic in Korea. The bulbs can be used similarly to garlic cloves, and the flowers are also edible. Parts of the plant can be used for preparing Van herbed cheese, a speciality of the Van province in Turkey.
Popular dishes using the plant include pesto, soups, pasta, cheese, scones and Devonnaise.
The leaves are also used as fodder. Cows that have fed on ramsons give milk that tastes slightly of garlic, and butter made from this milk used to be very popular in 19th-century Switzerland.
The first evidence of the human use of A. ursinum comes from the Mesolithic settlement of Barkær (Denmark), where an impression of a leaf has been found. In the Swiss Neolithic settlement of Thayngen-Weier (Cortaillod culture), a high concentration of pollen from A. ursinum was found in the settlement layer, interpreted by some as evidence for use of the plant as fodder.
Herbal remedy
Allium ursinum has been credited with many medicinal qualities and is a popular homeopathic ingredient. It is often used for treating cardiovascular, respiratory, and digestive problems, as well as for the sterilisation of wounds.
Various minerals are found in much higher amounts in Allium ursinum than in clove garlic. It is sometimes called the "magnesium king" of plants because of the high levels of this mineral found in the leaves.
Gallery
| Biology and health sciences | Asparagales | Plants |
477292 | https://en.wikipedia.org/wiki/Copper%28II%29%20sulfate | Copper(II) sulfate | Copper(II) sulfate is an inorganic compound with the chemical formula . It forms hydrates , where n can range from 1 to 7. The pentahydrate (n = 5), a bright blue crystal, is the most commonly encountered hydrate of copper(II) sulfate, while its anhydrous form is white. Older names for the pentahydrate include blue vitriol, bluestone, vitriol of copper, and Roman vitriol. It exothermically dissolves in water to give the aquo complex , which has octahedral molecular geometry. The structure of the solid pentahydrate reveals a polymeric structure wherein copper is again octahedral but bound to four water ligands. The centers are interconnected by sulfate anions to form chains.
Preparation and occurrence
Copper sulfate is produced industrially by treating copper metal with hot concentrated sulfuric acid or copper oxides with dilute sulfuric acid. For laboratory use, copper sulfate is usually purchased. Copper sulfate can also be produced by slowly leaching low-grade copper ore in air; bacteria may be used to hasten the process.
Commercial copper sulfate is usually about 98% pure copper sulfate, and may contain traces of water. Anhydrous copper sulfate is 39.81% copper and 60.19% sulfate by mass, and in its blue, hydrous form, it is 25.47% copper, 38.47% sulfate (12.82% sulfur) and 36.06% water by mass. Four types of crystal size are provided based on its usage: large crystals (10–40 mm), small crystals (2–10 mm), snow crystals (less than 2 mm), and windswept powder (less than 0.15 mm).
Chemical properties
Copper(II) sulfate pentahydrate decomposes before melting. It loses two water molecules upon heating at , followed by two more at and the final water molecule at .
The chemistry of aqueous copper sulfate is simply that of copper aquo complex, since the sulfate is not bound to copper in such solutions. Thus, such solutions react with concentrated hydrochloric acid to give tetrachlorocuprate(II):
Similarly treatment of such solutions with zinc gives metallic copper, as described by this simplified equation:
A further illustration of such single metal replacement reactions occurs when a piece of iron is submerged in a solution of copper sulfate:
In high school and general chemistry education, copper sulfate is used as an electrolyte for galvanic cells, usually as a cathode solution. For example, in a zinc/copper cell, copper ion in copper sulfate solution absorbs electron from zinc and forms metallic copper.
, E°cell = 0.34 V
Copper sulfate is commonly included in teenage chemistry sets and undergraduate experiments. It is often used to grow crystals in schools and in Copper electroplating experiments despite its toxicity. Copper sulfate is often used to demonstrate an exothermic reaction, in which steel wool or magnesium ribbon is placed in an aqueous solution of . It is used to demonstrate the principle of mineral hydration. The pentahydrate form, which is blue, is heated, turning the copper sulfate into the anhydrous form which is white, while the water that was present in the pentahydrate form evaporates. When water is then added to the anhydrous compound, it turns back into the pentahydrate form, regaining its blue color. Copper(II) sulfate pentahydrate can easily be produced by crystallization from solution as copper(II) sulfate, which is hygroscopic.
Uses
As a fungicide and herbicide
Copper sulfate has been used for control of algae in lakes and related fresh waters subject to eutrophication. It "remains the most effective algicidal treatment".
Bordeaux mixture, a suspension of copper(II) sulfate () and calcium hydroxide (), is used to control fungus on grapes, melons, and other berries. It is produced by mixing a water solution of copper sulfate and a suspension of slaked lime.
A dilute solution of copper sulfate is used to treat aquarium fishes for parasitic infections, and is also used to remove snails from aquariums and zebra mussels from water pipes. Copper ions are highly toxic to fish. Most species of algae can be controlled with very low concentrations of copper sulfate.
Analytical reagent
Several chemical tests utilize copper sulfate. It is used in Fehling's solution and Benedict's solution to test for reducing sugars, which reduce the soluble blue copper(II) sulfate to insoluble red copper(I) oxide. Copper(II) sulfate is also used in the Biuret reagent to test for proteins.
Copper sulfate is used to test blood for anemia. The blood is dropped into a solution of copper sulfate of known specific gravity—blood with sufficient hemoglobin sinks rapidly due to its density, whereas blood which sinks slowly or not at all has an insufficient amount of hemoglobin. Clinically relevant, however, modern laboratories utilize automated blood analyzers for accurate quantitative hemoglobin determinations, as opposed to older qualitative means.
In a flame test, the copper ions of copper sulfate emit a deep green light, a much deeper green than the flame test for barium.
Organic synthesis
Copper sulfate is employed at a limited level in organic synthesis. The anhydrous salt is used as a dehydrating agent for forming and manipulating acetal groups. The hydrated salt can be intimately mingled with potassium permanganate to give an oxidant for the conversion of primary alcohols.
Rayon production
Reaction with ammonium hydroxide yields tetraamminecopper(II) sulfate or Schweizer's reagent which was used to dissolve cellulose in the industrial production of Rayon.
Niche uses
Copper(II) sulfate has attracted many niche applications over the centuries. In industry copper sulfate has multiple applications. In printing it is an additive to book-binding pastes and glues to protect paper from insect bites; in building it is used as an additive to concrete to improve water resistance and prevent plant and mushroom growth. Copper sulfate can be used as a coloring ingredient in artworks, especially glasses and potteries. Copper sulfate is also used in firework manufacture as a blue coloring agent, but it is not safe to mix copper sulfate with chlorates when mixing firework powders.
Copper sulfate was once used to kill bromeliads, which serve as mosquito breeding sites. Copper sulfate is used as a molluscicide to treat bilharzia in tropical countries.
Art
In 2008, the artist Roger Hiorns filled an abandoned waterproofed council flat in London with 75,000 liters of copper(II) sulfate water solution. The solution was left to crystallize for several weeks before the flat was drained, leaving crystal-covered walls, floors and ceilings. The work is titled Seizure. Since 2011, it has been on exhibition at the Yorkshire Sculpture Park.
Etching
Copper(II) sulfate is used to etch zinc, aluminium, or copper plates for intaglio printmaking.
It is also used to etch designs into copper for jewelry, such as for Champlevé.
Dyeing
Copper(II) sulfate can be used as a mordant in vegetable dyeing. It often highlights the green tints of the specific dyes.
Electronics
An aqueous solution of copper(II) sulfate is often used as the resistive element in liquid resistors.
In electronic and microelectronic industry a bath of and sulfuric acid () is often used for electrodeposition of copper.
Other forms of copper sulfate
Anhydrous copper(II) sulfate can be produced by dehydration of the commonly available pentahydrate copper sulfate. In nature, it is found as the very rare mineral known as chalcocyanite. The pentahydrate also occurs in nature as chalcanthite. Other rare copper sulfate minerals include bonattite (trihydrate), boothite (heptahydrate), and the monohydrate compound poitevinite. There are numerous other, more complex, copper(II) sulfate minerals known, with environmentally important basic copper(II) sulfates like langite and posnjakite.
Toxicological effects
Copper(II) salts have an LD50 of 100 mg/kg.
Copper(II) sulfate was used in the past as an emetic. It is now considered too toxic for this use. It is still listed as an antidote in the World Health Organization's Anatomical Therapeutic Chemical Classification System.
| Physical sciences | Sulfuric oxyanions | Chemistry |
477661 | https://en.wikipedia.org/wiki/Sublimation%20%28phase%20transition%29 | Sublimation (phase transition) | Sublimation is the transition of a substance directly from the solid to the gas state, without passing through the liquid state. The verb form of sublimation is sublime, or less preferably, sublimate. Sublimate also refers to the product obtained by sublimation. The point at which sublimation occurs rapidly (for further details, see below) is called critical sublimation point, or simply sublimation point. Notable examples include sublimation of dry ice at room temperature and atmospheric pressure, and that of solid iodine with heating.
The reverse process of sublimation is deposition (also called desublimation), in which a substance passes directly from a gas to a solid phase, without passing through the liquid state.
Technically, all solids may sublime, though most sublime at extremely low rates that are hardly detectable under usual conditions. At normal pressures, most chemical compounds and elements possess three different states at different temperatures. In these cases, the transition from the solid to the gas state requires an intermediate liquid state. The pressure referred to is the partial pressure of the substance, not the total (e.g. atmospheric) pressure of the entire system. Thus, any solid can sublime if its vapour pressure is higher than the surrounding partial pressure of the same substance, and in some cases, sublimation occurs at an appreciable rate (e.g. water ice just below 0 °C).
For some substances, such as carbon and arsenic, sublimation from solid state is much more achievable than evaporation from liquid state and it is difficult to obtain them as liquids. This is because the pressure of their triple point in its phase diagram (which corresponds to the lowest pressure at which the substance can exist as a liquid) is very high.
Sublimation is caused by the absorption of heat which provides enough energy for some molecules to overcome the attractive forces of their neighbors and escape into the vapor phase. Since the process requires additional energy, sublimation is an endothermic change. The enthalpy of sublimation (also called heat of sublimation) can be calculated by adding the enthalpy of fusion and the enthalpy of vaporization.
Confusions
While the definition of sublimation is simple, there is often confusion as to what counts as a sublimation.
False correspondence with vaporization
Vaporization (from liquid to gas) is divided into two types: vaporization on the surface of the liquid is called evaporation, and vaporization at the boiling point with formation of bubbles in the interior of the liquid is called boiling. However there is no such distinction for the solid-to-gas transition, which is always called sublimation in both corresponding cases.
Potential distinction
For clarification, a distinction between the two corresponding cases is needed. With reference to a phase diagram, the sublimation that occurs left of the solid-gas boundary, the triple point or the solid-liquid boundary (corresponding to evaporation in vaporization) may be called gradual sublimation; and the substance sublimes gradually, regardless of rate. The sublimation that occurs at the solid-gas boundary (critical sublimation point) (corresponding to boiling in vaporization) may be called rapid sublimation, and the substance sublimes rapidly. The words "gradual" and "rapid" have acquired special meanings in this context and no longer describe the rate of sublimation.
Misuse for chemical reaction
The term sublimation refers specifically to a physical change of state and is not used to describe the transformation of a solid to a gas in a chemical reaction. For example, the dissociation on heating of solid ammonium chloride into hydrogen chloride and ammonia is not sublimation but a chemical reaction. Similarly the combustion of candles, containing paraffin wax, to carbon dioxide and water vapor is not sublimation but a chemical reaction with oxygen.
Historical definition
Sublimation is historically used as a generic term to describe a two-step phase transition ― a solid-to-gas transition (sublimation in a more precise definition) followed by a gas-to-solid transition (deposition). (See below)
Examples
The examples shown are substances that noticeably sublime under certain conditions.
Carbon dioxide
Solid carbon dioxide (dry ice) sublimes rapidly along the solid-gas boundary (sublimation point) below the triple point (e.g., at the temperature of −78.5 °C, at atmospheric pressure), whereas its melting into liquid CO2 can occur along the solid-liquid boundary (melting point) at pressures and temperatures above the triple point (i.e., 5.1 atm, −56.6 °C).
Water
Snow and ice sublime gradually at temperatures below the solid-liquid boundary (melting point) (generally 0 °C), and at partial pressures below the triple point pressure of , at a low rate. In freeze-drying, the material to be dehydrated is frozen and its water is allowed to sublime under reduced pressure or vacuum. The loss of snow from a snowfield during a cold spell is often caused by sunshine acting directly on the upper layers of the snow. Sublimation of ice is a factor to the erosive wear of glacier ice, known as ablation in glaciology.
Naphthalene
Naphthalene, an organic compound commonly found in pesticides such as mothballs, sublimes easily because it is made of non-polar molecules that are held together only by van der Waals intermolecular forces. Naphthalene is a solid that sublimes gradually at standard temperature and pressure, at a high rate, with the critical sublimation point at around . At low temperature, its vapour pressure is high enough, 1mmHg at 53°C, to make the solid form of naphthalene evaporate into gas. On cool surfaces, the naphthalene vapours will solidify to form needle-like crystals.
Iodine
Iodine sublimes gradually and produces visible fumes on gentle heating at standard atmospheric temperature. It is possible to obtain liquid iodine at atmospheric pressure by controlling the temperature at just between the melting point and the boiling point of iodine. In forensic science, iodine vapor can reveal latent fingerprints on paper.
Other substances
At atmospheric pressure, arsenic sublimes gradually upon heating, and sublimes rapidly at .
Cadmium and zinc sublime much more than other common materials, so they are not suitable materials for use in vacuum.
Purification by sublimation
Sublimation is a technique used by chemists to purify compounds. A solid is typically placed in a sublimation apparatus and heated under vacuum. Under this reduced pressure, the solid volatilizes and condenses as a purified compound on a cooled surface (cold finger), leaving a non-volatile residue of impurities behind. Once heating ceases and the vacuum is removed, the purified compound may be collected from the cooling surface.
For even higher purification efficiencies, a temperature gradient is applied, which also allows for the separation of different fractions. Typical setups use an evacuated glass tube that is heated gradually in a controlled manner. The material flow is from the hot end, where the initial material is placed, to the cold end that is connected to a pump stand. By controlling temperatures along the length of the tube, the operator can control the zones of re-condensation, with very volatile compounds being pumped out of the system completely (or caught by a separate cold trap), moderately volatile compounds re-condensing along the tube according to their different volatilities, and non-volatile compounds remaining in the hot end.
Vacuum sublimation of this type is also the method of choice for purification of organic compounds for use in the organic electronics industry, where very high purities (often > 99.99%) are needed to satisfy the standards for consumer electronics and other applications.
Historical usage
In ancient alchemy, a protoscience that contributed to the development of modern chemistry and medicine, alchemists developed a structure of basic laboratory techniques, theory, terminology, and experimental methods. Sublimation was used to refer to the process in which a substance is heated to a vapor, then immediately collects as sediment on the upper portion and neck of the heating medium (typically a retort or alembic), but can also be used to describe other similar non-laboratory transitions. It was mentioned by alchemical authors such as Basil Valentine and George Ripley, and in the Rosarium philosophorum, as a process necessary for the completion of the magnum opus. Here, the word sublimation was used to describe an exchange of "bodies" and "spirits" similar to laboratory phase transition between solids and gases. Valentine, in his Le char triomphal de l'antimoine (Triumphal Chariot of Antimony, published 1646) made a comparison to spagyrics in which a vegetable sublimation can be used to separate the spirits in wine and beer. Ripley used language more indicative of the mystical implications of sublimation, indicating that the process has a double aspect in the spiritualization of the body and the corporalizing of the spirit. He writes:
And Sublimations we make for three causes,
The first cause is to make the body spiritual.
The second is that the spirit may be corporeal,
And become fixed with it and consubstantial.
The third cause is that from its filthy original.
It may be cleansed, and its saltiness sulphurious
May be diminished in it, which is infectious.
Sublimation predictions
The enthalpy of sublimation has commonly been predicted using the equipartition theorem. If the lattice energy is assumed to be approximately half the packing energy, then the following thermodynamic corrections can be applied to predict the enthalpy of sublimation. Assuming a 1 molar ideal gas gives a correction for the thermodynamic environment (pressure and volume) in which pV = RT, hence a correction of 1RT. Additional corrections for the vibrations, rotations and translation then need to be applied. From the equipartition theorem gaseous rotation and translation contribute 1.5RT each to the final state, therefore a +3RT correction. Crystalline vibrations and rotations contribute 3RT each to the initial state, hence −6RT. Summing the RT corrections; −6RT + 3RT + RT = −2RT. This leads to the following approximate sublimation enthalpy. A similar approximation can be found for the entropy term if rigid bodies are assumed.
Dye-sublimation printing
Dye-sub printing is a digital printing technology using full color artwork that works with polyester and polymer-coated substrates. Also referred to as digital sublimation, the process is commonly used for decorating apparel, signs and banners, as well as novelty items such as cell phone covers, plaques, coffee mugs, and other items with sublimation-friendly surfaces. The process uses the science of sublimation, in which heat and pressure are applied to a solid, turning it into a gas through an endothermic reaction without passing through the liquid phase.
In sublimation printing, unique sublimation dyes are transferred to sheets of “transfer” paper via liquid gel ink through a piezoelectric print head. The ink is deposited on these high-release inkjet papers, which are used for the next step of the sublimation printing process. After the digital design is printed onto sublimation transfer sheets, it is placed on a heat press along with the substrate to be sublimated.
In order to transfer the image from the paper to the substrate, it requires a heat press process that is a combination of time, temperature and pressure. The heat press applies this special combination, which can change depending on the substrate, to “transfer” the sublimation dyes at the molecular level into the substrate. The most common dyes used for sublimation activate at 350 degrees Fahrenheit. However, a range of 380 to 420 degrees Fahrenheit is normally recommended for optimal color.
The result of the sublimation process is a nearly permanent, high resolution, full color print. Because the dyes are infused into the substrate at the molecular level, rather than applied at a topical level (such as with screen printing and direct to garment printing), the prints will not crack, fade or peel from the substrate under normal conditions.
| Physical sciences | Phase transitions | Physics |
477730 | https://en.wikipedia.org/wiki/Proboscis%20monkey | Proboscis monkey | The proboscis monkey or long-nosed monkey (Nasalis larvatus) is an arboreal Old World monkey with an unusually large nose, a reddish-brown skin color and a long tail. It is endemic to the southeast Asian island of Borneo and is found mostly in mangrove forests and on the coastal areas of the island.
This species co-exists with the Bornean orangutan and monkeys such as the silvery lutung. It belongs in the monotypic genus Nasalis.
Taxonomy
The proboscis monkey belongs to the subfamily Colobinae of the Old World monkeys. The two subspecies are:
N. l. larvatus (Wurmb, 1787), which occupies the whole range of the species excluding northeast Kalimantan
N. l. orientalis (Chasen, 1940), restricted to north-east Kalimantan
However, the difference between the subspecies is small, and not all authorities recognise N. l. orientalis.
The genus name Nasalis comes from the Latin word nasus meaning "nose".
This animal was made known to Westerners by Baron Friedrich von Wurmb in 1781, he later sent specimens of it to Stamford Raffles in Europe.
Description
The proboscis monkey is a large species, being one of the largest monkey species native to Asia. Only the Tibetan macaque and a few of the gray langurs can rival its size.
Sexual dimorphism is pronounced in the species. Males have a head-body length of and typically weigh , with a maximum known weight of . Females measure in head-and-body length and weigh , with a maximum known mass of . The male has a red penis with a black scrotum.
The proboscis monkey has a long coat; the fur on the back is bright orange, reddish brown, yellowish brown or brick-red. The underfur is light-grey, yellowish, or greyish to light-orange. Infants are born with a blue coloured face that at 2.5 months darkens to grey. By 8.5 months of age, the face has become cream coloured like the adults. Both sexes have bulging stomachs that give the monkeys what resembles a pot belly. Many of the monkeys' toes are webbed.
Nose
Further adding to the dimorphism is the large nose or proboscis of the male, which can exceed in length, and hangs lower than the mouth. Theories for the extensive length of their nose suggest it may be sexual selection by the females, who prefer louder vocalisations, with the size of the nose increasing the volume of the call.
The nose is smaller in the female and is upturned in the young. Nevertheless, the nose of the female is still fairly large for a primate. The skull of the proboscis monkey has specialized nasal cartilages that support the large nose.
Distribution and habitat
The proboscis monkey is endemic to the island of Borneo and can be found in all three nations that divide the island: Brunei, Indonesia and Malaysia. It is most common in coastal areas and along rivers. This species is restricted to lowland habitats that may experience tides. It favors dipterocarp, mangrove and riverine forests. It can also be found in swamp forests, stunted swamp forests, rubber forests, rubber plantations, limestone hill forests, nypa swamps, nibong swamps, and tall swamp forests, tropical heath forests and steep cliffs. This species usually stays within at least a kilometer from a water source. It is perhaps the most aquatic of the primates and is a fairly good swimmer, capable of swimming up to underwater. It is known to swim across rivers. Aside from this, the proboscis monkey is largely arboreal and moves quadrupedally and by leaps. It is known to jump off branches and descend into water.
Behavior and ecology
Social behavior
Proboscis monkeys generally live in groups composed of one adult male, some adult females and their offspring. All-male groups may also exist. Some individuals are solitary, mostly males. Monkey groups live in overlapping home ranges, with little territoriality, in a fission-fusion society, with groups gathering at sleeping sites as night falls. There exist bands which arise when groups come together and slip apart yet sometimes groups may join to mate and groom. Groups gather during the day and travel together, but individuals only groom and play with those in their own group. One-male groups consist of 3 to 19 individuals, while bands can consist of as many as 60 individuals. Serious aggression is uncommon among the monkeys but minor aggression does occur. Overall, members of the same bands are fairly tolerant of each other. A linear dominance hierarchy exists between females. Males of one-male groups can stay in their groups for six to eight years. Replacements in the resident males appear to occur without serious aggression. Upon reaching adulthood, males leave their natal groups and join all-male groups. Females also sometimes leave their natal groups, perhaps to avoid infanticide or inbreeding, reduce competition for food, or elevation of their social status.
In Sabah, Malaysia, proboscis monkeys have been observed in mixed-species groups with silvery lutungs, and interspecific mating and a possible hybrid has been observed. Researchers believe this may be a result of the two species being confined to a small patch of riverine forest due to deforestation in order to plant oil palm trees.
Reproduction
Females become sexually mature at the age of five years. They experience sexual swelling, which involves the genitals becoming pink or reddened. At one site, matings largely take place between February and November, while births occur between March and May. Copulations tend to last for half a minute. The male will grab the female by the ankles or torso and mount her from behind. Both sexes will encourage mating, but they are not always successful. When soliciting, both sexes will make pouted faces. In addition, males will sometimes vocalize and females will present their backsides and shake their head from side to side. Mating pairs are sometimes harassed by subadults. Proboscis monkeys may also engage in mounting with no reproductive purpose, such as playful and same-sex mounting, and females will attempt to initiate copulation even after they have conceived. Gestation usually last 166–200 days or slightly more. Females tend to give birth at night or in the early morning. The mothers then eat the placenta and lick their infants clean. The young begin to eat solid foods at six weeks and are weaned at seven months old. The nose of a young male grows slowly until reaching adulthood. The mother will allow other members of her group to hold her infant. When a resident male in a one-male group is replaced, the infants are at risk of infanticide.
Communication
Proboscis monkeys are known to make various vocalizations. When communicating the status of group, males will emit honks. They have a special honk emitted towards infants, which is also used for reassurance. Males will also produce alarm calls to signal danger. Both sexes give threat calls, but each are different. In addition, females and immature individuals will emit so-called "female calls" when angry. Honks, roars and snarls are made during low-intensity agonistic encounters. Nonvocal displays include leaping-branch shaking, bare-teeth open mouth threats and erection in males, made in the same situations.
Feeding and activities
As a seasonal folivore and frugivore, the proboscis monkey eats primarily fruit and leaves. It also eats flowers, seeds and insects to a lesser extent. At least 55 different plant species are consumed, "with a marked preference for Eugenia sp., Ganua motleyana and Lophopetalum javanicum". Young leaves are preferred over mature leaves and unripe fruits are preferred over ripe fruit. Being a seasonal eater, the proboscis monkey eats mostly fruit from January to May and mostly leaves from June to December. Groups usually sleep in adjacent trees. Monkeys tend to sleep near rivers, if they are nearby. Proboscis monkeys will start the day foraging and then rest further inland. Their daily activities consist of resting, traveling, feeding and keeping vigilant. Occasionally, they chew their cud to allow more efficient digestion and food intake. As night approaches, the monkeys move back near the river and forage again. Predators (potential or confirmed) of the proboscis monkey include crocodilians like false gharials and saltwater crocodiles, the Sunda clouded leopard, sun bears and reticulated pythons as well as, for probably young or sickly monkeys, large eagles (such as the crested serpent eagle or black eagle), large owls, and monitor lizards. Monkeys will cross rivers at narrows or cross arboreally if possible. This may serve as predator avoidance.
Conservation status
The proboscis monkey is assessed as endangered in the IUCN Red List of Threatened Species and listed in Appendix I of CITES. Its total population has decreased by more than 50% in the past 36–40 years to 2008 due to ongoing habitat loss because of logging and oil palm plantations, and hunting in some areas due to the species being treated as a delicacy, as well as its use in traditional Chinese medicine. The population is fragmented: the largest remaining populations are found in Kalimantan; there are far fewer in Sarawak, Brunei and Sabah. The proboscis monkey is protected by law in all regions of Borneo. In Malaysia, it is protected by a number of laws including the Wildlife Protection Act (federal law), the Wildlife Protection Ordinance 1998 (Chapter 26) and Wildlife Conservation Enactment 1997 (Sabah state law).
The proboscis monkey can be found in 16 protected areas: Danau Sentarum National Park, Gunung Palung National Park, Kendawangan Nature Reserve, Kutai National Park, Lesan Protection Forest, Muara Kaman Nature Reserve, Mandor Reserve and Tanjung Puting National Park in Indonesia; Bako National Park, Gunung Pueh Forest Reserve, Kabili-Sepilok Forest Reserve, Klias National Park, Kulamba Wildlife Reserve, Lower Kinabatangan Wildlife Sanctuary, Sungei Samunsam Wildlife Sanctuary and Ulu Segama Reserve in Malaysia.
| Biology and health sciences | Old World monkeys | Animals |
2607068 | https://en.wikipedia.org/wiki/Giant%20otter | Giant otter | The giant otter or giant river otter (Pteronura brasiliensis) is a South American carnivorous mammal. It is the longest member of the weasel family, Mustelidae, a globally successful group of predators, reaching up to . Atypical of mustelids, the giant otter is a social species, with family groups typically supporting three to eight members. The groups are centered on a dominant breeding pair and are extremely cohesive and cooperative. Although generally peaceful, the species is territorial, and aggression has been observed between groups. The giant otter is diurnal, being active exclusively during daylight hours. It is the noisiest otter species, and distinct vocalizations have been documented that indicate alarm, aggression, and reassurance.
The giant otter ranges across north-central South America; it lives mostly in and along the Amazon River and in the Pantanal. Its distribution has been greatly reduced and is now discontinuous. Decades of poaching for its velvety pelt, peaking in the 1950s and 1960s, considerably diminished population numbers. The species was listed as endangered in 1999 and wild population estimates are typically below 5,000. The Guianas are one of the last real strongholds for the species, which also enjoys modest numbers – and significant protection – in the Peruvian Amazonian basin. It is one of the most endangered mammal species in the Neotropics. Habitat degradation and loss is the greatest current threat. They are also rare in captivity; in 2003, only 60 giant otters were being held.
The giant otter shows a variety of adaptations suitable to an amphibious lifestyle, including exceptionally dense fur, a wing-like tail, and webbed feet. The species prefers freshwater rivers and streams, which are usually seasonally flooded, and may also take to freshwater lakes and springs. It constructs extensive campsites close to feeding areas, clearing large amounts of vegetation. The giant otter subsists almost exclusively on a diet of fish, particularly characins and catfish, but may also eat crabs, turtles, snakes and small caimans. It has no serious natural predators other than humans, although it must compete with other predators, such as the Neotropical otters and various crocodilian species, for food resources.
Name
The giant otter has a handful of other names. In Brazil it is known as ariranha, from the Tupi word , or onça-d'água, meaning water jaguar. In Spanish, river wolf () and water dog () are used occasionally (though the latter also refers to several different animals) and may have been more common in the reports of explorers in the 19th and early 20th centuries. All four names are in use in South America, with a number of regional variations. "Giant otter" translates literally as and in Spanish and Portuguese, respectively. Among the Achuar people, they are known as wankanim, among the Sanumá as hadami, and among the Makushi as turara. The genus name, Pteronura, is derived from the Ancient Greek words (, feather or wing) and (, tail), a reference to its distinctive, wing-like tail.
Taxonomy and evolution
The otters form the subfamily Lutrinae within the mustelids and the giant otter is the only member of the genus Pteronura. Two subspecies are currently recognized by the canonical Mammal Species of the World, P. b. brasiliensis and P. b. paraguensis. Incorrect descriptions of the species have led to multiple synonyms (the latter subspecies is often P. b. paranensis in the literature). P. b. brasiliensis is distributed across the north of the giant otter range, including the Orinoco, Amazon, and Guianas river systems; to the south, P. b. paraguensis has been suggested in Paraguay, Uruguay, southern Brazil, and northern Argentina, although it may be extinct in the last three of these four. The International Union for Conservation of Nature (IUCN) considers the species' presence in Argentina and Uruguay uncertain. In the former, investigation has shown thinly distributed population remnants. P. b. paraguensis is supposedly smaller and more gregarious, with different dentition and skull morphology. Carter and Rosas, however, rejected the subspecific division in 1997, noting the classification had only been validated once, in 1968, and the P. b. paraguensis type specimen was very similar to P. b. brasiliensis. Biologist Nicole Duplaix calls the division of "doubtful value".
The earliest fossil evidence of the giant river otter dates to the Late Pleistocene of Argentina, and it was slightly larger than known modern specimens. An extinct genus, Satherium, is believed to be ancestral to the present species, having migrated to the New World during the Pliocene or early Pleistocene. The giant otter shares the South American continent with three of the four members of the New World otter genus Lontra: the Neotropical river otter, the southern river otter, and the marine otter. (The North American river otter (Lontra canadensis) is the fourth Lontra member.) The giant otter seems to have evolved independently of Lontra in South America, despite the overlap. The smooth-coated otter (Lutrogale perspicillata) of Asia may be its closest extant relative; similar behaviour, vocalizations, and skull morphology have been noted. Both species also show strong pair bonding and paternal engagement in rearing cubs. Giant otter fossil remains have been recovered from a cave in the Brazilian Mato Grosso.
Phylogenetic analysis by Koepfli and Wayne in 1998 found the giant otter has the highest divergence sequences within the otter subfamily, forming a distinct clade that split away 10 to 14 million years ago. They noted that the species may be the basal divergence among the otters or fall outside of them altogether, having split even before other mustelids, such as the ermine, polecat, and mink. Later gene sequencing research on the mustelids, from 2004, places the divergence of the giant otter somewhat later, between five and 11 million years ago; the corresponding phylogenetic tree locates the Lontra divergence first among otter genera, and Pteronura second, although divergence ranges overlap.
Physical characteristics
The giant otter is clearly distinguished from other otters by morphological and behavioural characteristics. It has the greatest body length of any species in the mustelid family, although the sea otter may be heavier. Males are between in length from head to tail and females between . The animal's well-muscled tail can add a further to the total body length. Early reports of skins and living animals suggested exceptionally large males of up to ; intensive hunting likely reduced the occurrence of such massive specimens. Weights are between for males and for females.
The giant otter has the shortest fur of all otter species; it is typically chocolate brown, but may be reddish or fawn, and appears nearly black when wet. The fur is extremely dense, so much so that water cannot penetrate to the skin. Guard hairs trap water and keep the inner fur dry; the guard hairs are approximately 8 millimetres (one-third of an inch) in length, about twice as long as the fur of the inner coat. Its velvety feel makes the animal highly sought after by fur traders and has contributed to its decline. Individuals can be identified from birth by unique markings of white or cream fur on the throat and under the chin.
Giant otter muzzles are short and sloping and give the head a ball-shaped appearance. The ears are small and rounded. The nose (or rhinarium) is completely covered in fur, with only the two slit-like nostrils visible. The giant otter's highly sensitive whiskers (vibrissae) allow the animal to track changes in water pressure and currents, which aids in detecting prey. The legs are short and stubby and end in large webbed feet tipped with sharp claws. Well suited for an aquatic life, it can close its ears and nose while underwater.
At the time of Carter and Rosas's writing, vision had not been directly studied, but field observations show the animal primarily hunts by sight; above water, it is able to recognize observers at great distances. The fact that it is exclusively active during the day further suggests its eyesight should be strong, to aid in hunting and predator avoidance. In other otter species, vision is generally normal or slightly myopic, both on land and in water. The giant otter's hearing is acute and its sense of smell is excellent.
The species possesses 2n = 38 chromosomes.
Biology and behaviour
The giant otter is large, gregarious, and diurnal. Early travelers' reports describe noisy groups surrounding explorers' boats, but little scientific information was available on the species until Duplaix's groundbreaking work in the late 1970s. Concern over this endangered species has since generated a body of research.
Vocalizations
The giant otter is an especially noisy animal, with a complex repertoire of vocalizations. All otters produce vocalizations, but by frequency and volume, the giant otter may be the most vocal. Duplaix identified nine distinct sounds, with further subdivisions possible, depending on context. Quick hah barks or explosive snorts suggest immediate interest and possible danger. A wavering scream may be used in bluff charges against intruders, while a low growl is used for aggressive warning. Hums and coos are more reassuring within the group. Whistles may be used as advance warning of nonhostile intent between groups, although evidence is limited. Newborn pups squeak to elicit attention, while older young whine and wail when they begin to participate in group activities. An analysis published in 2014 cataloged 22 distinct types of vocalization in adults and 11 in neonates. Each family of otters was shown to have its own unique audio signature.
Social structure
The giant otter is a highly social animal and lives in extended family groups. Group sizes are anywhere from two to 20 members, but likely average between three and eight. (Larger figures may reflect two or three family groups temporarily feeding together.)
Group members share roles, structured around the dominant breeding pair. The species is territorial, with groups marking their ranges with latrines, gland secretions, and vocalizations. At least one case of a change in alpha relationship has been reported, with a new male taking over the role; the mechanics of the transition were not determined. Duplaix suggests a division between "residents", who are established within groups and territories, and nomadic and solitary "transients"; the categories do not seem rigid, and both may be a normal part of the giant otter life cycle. One tentative theory for the development of sociality in mustelids is that locally abundant but unpredictably dispersed prey causes groups to form.
Aggression within the species ("intraspecific" conflict) has been documented. Defence against intruding animals appears to be cooperative: while adult males typically lead in aggressive encounters, cases of alpha females guarding groups have been reported. One fight was directly observed in the Brazilian Pantanal in which three animals violently engaged a single individual near a range boundary. In another instance in Brazil, a carcass was found with clear indications of violent assault by other otters, including bites to the snout and genitals, an attack pattern similar to that exhibited by captive animals. While not rare among large predators in general, intraspecific aggression is uncommon among otter species; Ribas and Mourão suggest a correlation to the animal's sociability, which is also rare among other otters.
A capacity for aggressive behavior should not be overstated with the giant otter. Researchers emphasize that even between groups, conflict avoidance is generally adopted. Within groups, the animals are extremely peaceful and cooperative. Group hierarchies are not rigid and the animals easily share roles.
Reproduction and life cycle
Giant otters build dens, which are holes dug into riverbanks, usually with multiple entrances and multiple chambers inside. They give birth within these dens during the dry season. In Cantão State Park, otters dig their reproductive dens on the shores of oxbow lakes starting around July, when waters are already quite low. They give birth between August and September, and the young pups emerge for the first time in October and November, which are the months of lowest water when fish concentrations in the dwindling lakes and channels are at their peak. This makes it easier for the adults to catch enough fish for the growing young, and for the pups to learn how to catch fish. The entire group, including nonreproductive adults, which are usually older siblings to that year's pups, collaborates to catch enough fish for the young.
Details of giant otter reproduction and life cycle are scarce, and captive animals have provided much of the information. Females appear to give birth year round, although in the wild, births may peak during the dry season. The estrous cycle is 21 days, with females receptive to sexual advances between three and 10 days. Study of captive specimens has found only males initiate copulation. At Tierpark Hagenbeck in Germany, long-term pair bonding and individualized mate selection were seen, with copulation most frequently taking place in water. Females have a gestation period of 65 to 70 days, giving birth to one to five pups, with an average of two. Research over five years on a breeding pair at the Cali Zoo in Colombia found the average interval between litters was six to seven months, but as short as 77 days when the previous litter did not survive. Other sources have found greater intervals, with as long as 21 to 33 months suggested for otters in the wild.
Mothers give birth to furred and blind cubs in an underground den near the river shore and fishing sites. Males actively participate in rearing cubs and family cohesion is strong; older, juvenile siblings also participate in rearing, although in the weeks immediately after birth, they may temporarily leave the group. Pups open their eyes in their fourth week, begin walking in their fifth, and are able to swim confidently between 12 and 14 weeks old. They are weaned by nine months and begin hunting successfully soon after. The animal reaches sexual maturity at about two years of age and both male and female pups leave the group permanently after two to three years. They then search for new territory to begin a family of their own.
Studies of giant otters in captivity have given indications about the environment necessary to both maintain a physically and behaviorally healthy population and allow successful cub-rearing. These include providing at least the minimum recommended land-to-water area ratio, and that all enclosure land surfaces (both artificial and natural) are nearly entirely covered with the recommended substrate conditions (e.g. tree-bark mulch and soft pebble-free sand/soil). Ensuring that the animals have sufficient privacy from human disturbances (visual and acoustic, from zoo staff or visitors) at parturition and during cub-rearing is also essential, but not sufficient. Insufficient land area proportions and unsuitable substrate conditions in zoos have historically been the primary cause of high cub mortality and physical and behavioral health problems among giant otters. For example, stress to the parents during cub-rearing due to inappropriate enclosure conditions has been the primary reason for cub neglect, abuse and infanticide.
In the wild, it has been suggested, although not systematically confirmed, that tourists cause similar stresses: disrupted lactation and denning, reduced hunting, and habitat abandonment are all risks. This sensitivity is matched by a strong protectiveness towards the young. All group members may aggressively charge intruders, including boats with humans in them.
The longest documented giant otter lifespan in the wild is eight years. In captivity, this may increase to 17, with an unconfirmed record of 19. The animal is susceptible to a variety of diseases, including canine parvovirus. Parasites, such as the larvae of flies and a variety of intestinal worms, also afflict the giant otter. Other causes of death include accidents, gastroenteritis, infanticide, and epileptic seizures.
Hunting and diet
The giant otter is an apex predator, and its population status reflects the overall health of riverine ecosystems. It feeds mainly on fish, including cichlids, perch, characins (such as piranha), and catfish. One full-year study of giant otter scats in Amazonian Brazil found fish present in all fecal samples. Fish from the order Perciformes, particularly cichlids and perch, were seen in 97% of scats, and Characiformes, such as characins, in 86%. Fish remains were of medium-sized species that seem to prefer relatively shallow water, to the advantage of the visually oriented giant otter. Prey species found were also sedentary, generally swimming only short distances, which may aid the giant otter in predation. Hunting in shallow water has also been found to be more rewarding, with water depth less than having the highest success rate. The giant otter seems to be opportunistic, taking whatever species are most locally abundant. If fish are unavailable, it will also take crabs, snakes, and even small caimans and anacondas.
The species can hunt singly, in pairs, and in groups, relying on sharp eyesight to locate prey. In some cases, supposed cooperative hunting may be incidental, a result of group members fishing individually in close proximity; truly coordinated hunting may only occur where the prey cannot be taken by a single giant otter, such as with small anacondas and juvenile black caiman. The giant otter seems to prefer prey fish that are generally immobile on river bottoms in clear water. Prey chase is rapid and tumultuous, with lunges and twists through the shallows and few missed targets. The otter can attack from both above and below, swiveling at the last instant to clamp the prey in its jaws. Giant otters catch their own food and consume it immediately; they grasp the fish firmly between the forepaws and begin eating noisily at the head. Carter and Rosas have found captive adult animals consume around 10% of their body weight daily—about , in keeping with findings in the wild.
Ecology
Habitat
The species is amphibious, although primarily terrestrial. It occurs in freshwater rivers and streams, which generally flood seasonally. Other water habitats include freshwater springs and permanent freshwater lakes. Four specific vegetation types occur on one important creek in Suriname: riverbank high forest, floodable mixed marsh and high swamp forest, floodable low marsh forest, and grass islands and floating meadows within open areas of the creek itself. Duplaix identified two critical factors in habitat selection: food abundance, which appears to positively correlate to shallow water, and low sloping banks with good cover and easy access to preferred water types. The giant otter seems to choose clear, black waters with rocky or sandy bottoms over silty, saline, and white waters.
Giant otters use areas beside rivers for building dens, campsites, and latrines. They clear significant amounts of vegetation while building their campsites. One report suggests maximum areas long and wide, well-marked by scent glands, urine, and feces to signal territory. Carter and Rosas found average areas a third this size. Giant otters adopt communal latrines beside campsites, and dig dens with a handful of entrances, typically under root systems or fallen trees. One report found between three and eight campsites, clustered around feeding areas. In seasonally flooded areas, the giant otter may abandon campsites during the wet season, dispersing to flooded forests in search of prey. Giant otters may adopt preferred locations perennially, often on high ground. These can become quite extensive, including "backdoor" exits into forests and swamps, away from the water. Otters do not visit or mark every site daily, but usually patrol all of them, often by a pair of otters in the morning.
Research generally takes place in the dry season and an understanding of the species' overall habitat use remains partial. An analysis of dry season range size for three otter groups in Ecuador found areas between . Utreras presumed habitat requirements and availability would differ dramatically in the rainy season: estimating range sizes of 1.98 to as much as 19.55 square kilometres (0.76 to 7.55 sq miles) for the groups. Other researchers suggest approximately and note a strong inverse correlation between sociality and home range size; the highly social giant otter has smaller home range sizes than would be expected for a species of its mass. Population densities varied with a high of reported in Suriname and with a low of found in Guyana.
In 2021, conservationists at Fundación Rewilding spotted a wild giant otter swimming in the Bermejo River in Impenetrable National Park, located in the Chaco province of northeast Argentina.
Predation and competition
Adult giant otters living in family groups have no known serious natural predators; however, there are some accounts of black caimans in Peru and yacare caimans in the Pantanal preying on giant otters. In addition, solitary animals and young may be vulnerable to attacks by the jaguar, cougar, and anaconda, but this is based on historical reports, not direct observation. Pups are more vulnerable, and may be taken by caiman and other large predators, although adults are constantly mindful of stray young, and will harass and fight off possible predators. When in the water, the giant otter faces danger from animals not strictly preying upon it: the electric eel and stingray are potentially deadly if stumbled upon, and piranha may be capable of at least taking bites out of a giant otter, as evidenced by scarring on individuals.
Even if without direct predation, the giant otter must still compete with other predators for food resources. Duplaix documented interaction with the Neotropical otter. While the two species are sympatric (with overlapping ranges) during certain seasons, there appeared to be no serious conflict. The smaller neotropical otter is far more shy, less noisy, and less social; at about a third the weight of the giant otter, it is more vulnerable to predation, hence, a lack of conspicuousness is to its advantage. The neotropical otter is active during twilight and darkness, reducing the likelihood of conflict with the diurnal giant otter. Its smaller prey, different denning habits, and different preferred water types also reduce interaction.
Other species that prey upon similar food resources include the caimans and large fish that are themselves piscivores. Gymnotids, such as the electric eel, and the large silurid catfish are among aquatic competitors. Two river dolphins, the tucuxi and Amazon river dolphin, might potentially compete with the giant otter, but different spatial use and dietary preferences suggest minimal overlap. Furthermore, Defler observed associations between giant otters and the Amazon river dolphins, and suggested that dolphins may benefit by fish fleeing from the otters. The spectacled caiman is another potential competitor, but Duplaix found no conflict with the species in Suriname.
Conservation status
The IUCN listed the giant otter as "endangered" in 1999; it had been considered "vulnerable" under all previous listings from 1982 when sufficient data had first become available. It is regulated internationally under Appendix I of the Convention on International Trade in Endangered Species of Wild Fauna and Flora (CITES) meaning commercial trade in specimens (including parts and derivatives) is prohibited.
Threats
The animal faces a variety of critical threats. Poaching has long been a problem. Statistics show between 1959 and 1969 Amazonian Brazil alone accounted for 1,000 to 3,000 pelts annually. The species was so thoroughly decimated, the number dropped to just 12 in 1971. The implementation of CITES in 1973 finally brought about significant hunting reductions, although demand did not disappear entirely: in the 1980s, pelt prices were as high as US$250 on the European market. The threat has been exacerbated by the otters' relative fearlessness and tendency to approach human beings. They are extremely easy to hunt, being active through the day and highly inquisitive. The animal's relatively late sexual maturity and complex social life makes hunting especially disastrous.
More recently, habitat destruction and degradation have become the principal dangers, and a further reduction of 50% is expected in giant otter numbers within the 25 years after 2020 (about the span of three generations of giant otters). Typically, loggers first move into rainforest, clearing the vegetation along riverbanks. Farmers follow, creating depleted soil and disrupted habitats. As human activity expands, giant otter home ranges become increasingly isolated. Subadults leaving in search of new territory find it impossible to set up family groups. Specific threats from human industry include unsustainable mahogany logging in parts of the giant otter range, and concentrations of mercury in its diet of fish, a byproduct of gold mining.
Other threats to the giant otter include conflict with fishermen, who often view the species as a nuisance (see below). Ecotourism also presents challenges: while it raises money and awareness for the animals, by its nature it also increases human effect on the species, both through associated development and direct disturbances in the field. A number of restrictions on land use and human intrusion are required to properly maintain wild populations. Schenck et al., who undertook extensive fieldwork in Peru in the 1990s, suggest specific "no-go" zones where the species is most frequently observed, offset by observation towers and platforms to allow viewing. Limits on the number of tourists at any one time, fishing prohibitions, and a minimum safe distance of are proposed to offer further protection.
Distribution and population
The giant otter has lost as much as 80% of its South American range. While still present in a number of north-central countries, giant otter populations are under considerable stress. The IUCN lists Bolivia, Brazil, Colombia, Ecuador, French Guiana, Guyana, Paraguay, Peru, Suriname, and Venezuela as current range countries. Given local extinctions, the species' range has become discontinuous. Total population numbers are difficult to estimate. Populations in Bolivia were once widespread but the country became a "black spot" on distribution maps after poaching between the 1940s and 1970s; a relatively healthy, but still small, population of 350 was estimated in the country in 2002. The species has likely been extirpated from southern Brazil, but in the west of the country, decreased hunting pressure in the critical Pantanal has led to very successful recolonization; an estimate suggests 1,000 or more animals in the region.
As of 2020, the IUCN estimates that there may be 4,569 otters living in Brazil. A significant population lives in the wetlands of the central Araguaia River, and in particular within Cantão State Park, which, with its 843 oxbow lakes and extensive flooded forests and marshlands, is one of the best habitat patches for this species in Brazil.
Suriname still has significant forest cover and an extensive system of protected areas, much of which protects the giant otter. Duplaix returned to the country in 2000 and found the giant otter still present on the Kaburi Creek, a "jewel" of biodiversity, although increased human presence and land use suggests, sooner or later, the species may not be able to find suitable habitat for campsites. In a report for World Wildlife Fund in 2002, Duplaix was emphatic about the importance of Suriname and the other Guianas:
Other countries have taken a lead in designating protected areas in South America. In 2004, Peru created one of the largest conservation areas in the world, Alto Purús National Park, with an area similar in size to Belgium. The park harbors many endangered plants and animals, including the giant otter, and holds the world record for mammal diversity. Bolivia designated wetlands larger than the size of Switzerland as a freshwater protected area in 2001; these are also home to the giant otter.
Interactions with indigenous peoples
Throughout its range, the giant otter interacts with indigenous groups, who often practice traditional hunting and fishing. A study of five indigenous communities in Colombia suggests native attitudes toward the animal are a threat: the otters are often viewed as a nuisance that interferes with fishing, and are sometimes killed. Even when told of the importance of the species to ecosystems and the danger of extinction, interviewees showed little interest in continuing to coexist with the species. Schoolchildren, however, had a more positive impression of the animal.
In Suriname, the giant otter is not a traditional prey species for human hunters, which affords some protection. (One researcher has suggested the giant otter is hunted only in desperation due to its horrible taste.) The animal sometimes drowns in nets set across rivers and machete attacks by fishermen have been noted, according to Duplaix, but "tolerance is the rule" in Suriname. One difference in behavior was seen in the country in 2002: the normally inquisitive giant otters showed "active avoidance behavior with visible panic" when boats appeared. Logging, hunting, and pup seizure may have led groups to be far more wary of human activity.
Local people sometimes take pups for the exotic pet trade or as pets for themselves, but the animal rapidly grows to become unmanageable. Duplaix relates the story of an Arawak Indian who took two pups from their parents. While revealing of the affection held for the animals, the seizure was a profound blow to the breeding pair, which went on to lose their territory to competitors.
The species has also appeared in the folklore of the region. It plays an important role in the mythology of the Achuar people, where giant otters are seen as a form of the tsunki, or water spirits: they are a sort of "water people" who feed on fish. They appear in a fish poisoning legend where they assist a man who has wasted his sexual energy, creating the anacondas of the world from his distressed and extended genitals.
The Bororó people have a legend on the origin of tobacco smoking: those who used the leaf improperly by swallowing it were punished by being transformed into giant otters; the Bororo also associate the giant otter with fish and with fire. A Ticuna legend has it that the giant otter exchanged places with the jaguar: the story says jaguar formerly lived in the water and the giant otter came to the land only to eat. The indigenous Kichwa peoples from Amazonian Peru believed in a world of water where Yaku runa reigned as mother of the water and was charged with caring for fish and animals. Giant otters served as Yaku runa's canoes. A Maxacali creation story suggests that the practice of otter fishing may have been prevalent in the past.
| Biology and health sciences | Mustelidae | Animals |
2607200 | https://en.wikipedia.org/wiki/Edaphosaurus | Edaphosaurus | Edaphosaurus (, meaning "pavement lizard" for dense clusters of teeth) is a genus of extinct edaphosaurid synapsids that lived in what is now North America and Europe around 303.4 to 272.5 million years ago, during the Late Carboniferous to Early Permian. American paleontologist Edward Drinker Cope first described Edaphosaurus in 1882, naming it for the "dental pavement" on both the upper and lower jaws, from the Greek ("ground"; also "pavement") and () ("lizard").
Edaphosaurus is important as one of the earliest-known, large, plant-eating (herbivorous), amniote tetrapods (four-legged land-living vertebrates). In addition to the large tooth plates in its jaws, the most characteristic feature of Edaphosaurus is a sail on its back. A number of other synapsids from the same time period also have tall dorsal sails, most famously the large apex predator Dimetrodon. However, the sail on Edaphosaurus is different in shape and morphology. The first fossils of Edaphosaurus came from the Texas Red Beds in North America, with later finds in New Mexico, Oklahoma, West Virginia, and Ohio. Fragmentary fossils attributed to Edaphosaurus have also been found in eastern Germany in Central Europe.
Etymology
The name Edaphosaurus, meant as "pavement lizard", is often translated inaccurately as "earth lizard", "ground lizard", or "foundation lizard" based on other meanings for the Greek , such as "soil, earth, ground, land, base" used in Neo-Latin scientific nomenclature (edaphology). However, older names in paleontology, such as Edaphodon Buckland, 1838 "pavement tooth" (a fossil fish), match Cope's clearly intended meaning "pavement" for Greek edaphos in reference to the animal's teeth.
Description and paleobiology
Edaphosaurus species measured from in length and weighed over 300 kg (660 lb). In keeping with its tiny head, the cervical vertebrae are reduced in length, while the dorsal vertebrae are massive, the tail is deep, the limbs are short and robust, and the ribs form a wide ribcage. Like most herbivores, Edaphosaurus would have had a capacious gut and symbiotic bacteria to aid in the breakdown of cellulose and other indigestible plant material. Like its more famous relative Dimetrodon, Edaphosaurus had a sail-like fin that was supported by bones of the vertebral column. Edaphosaurus differs from Dimetrodon in having cross-bars on the spines that supported its fin.
Skull
The head of Edaphosaurus was short, relatively broad, triangular in outline, and remarkably small compared to its body size. The deep lower jaw likely had powerful muscles and the marginal teeth along the front and sides of its jaws had serrated tips, helping Edaphosaurus to crop bite-sized pieces from tough terrestrial plants. Back parts of the roof of the mouth and the inside of the lower jaw held dense batteries of peglike teeth, forming a broad crushing and grinding surface on each side above and below. Its jaw movements were propalinal (front to back). Early descriptions suggested that Edaphosaurus fed on invertebrates such as mollusks, which it would have crushed with its tooth plates. However, paleontologists now think that Edaphosaurus ate plants, although tooth-on-tooth wear between its upper and lower tooth plates indicates only "limited processing of food" compared to other early plant-eaters such as Diadectes, a large nonamniote reptiliomorph (Diadectidae) that lived at the same time.
The recently described Melanedaphodon from the Middle Pennsylvanian subperiod of the Carboniferous Period in North America is currently the earliest known edaphosaurid and represents a transitional stage from a diet of hard-shelled invertebrates such as insects and mollusks to fibrous plants. Melanedaphodon had large and bulbous teeth along its upper and lower jaws, but also had "a moderately-developed tooth battery" on its palate, "which appears intermediary towards the condition seen in Edaphosaurus" and would have helped process tough plant material. Melanedaphodon was found to be a sister taxon to Edaphosaurus and lived earlier than the edaphosaurid Ianthasaurus, which lacked tooth plates and ate insects.
Sail
The sail along the back of Edaphosaurus was supported by hugely elongated neural spines from neck to lumbar region, connected by tissue in life. When compared with the sail of Dimetrodon, the vertebral spines are shorter and heavier, and bear numerous small crossbars. Edaphosaurus and other members of the Edaphosauridae evolved tall dorsal sails independently of sail-back members of the Sphenacodontidae such as Dimetrodon and Secodontosaurus that lived at the same time, an unusual example of parallel evolution. The of the sail in both groups is still debated. Researchers have suggested that such sails could have provided camouflage, wind-powered sailing over water, anchoring for extra muscle support and rigidity for the backbone, protection against predator attacks, fat-storage areas, body-temperature control surfaces, or sexual display and species recognition. The height of the sail, curvature of the spines, and shape of the crossbars are distinct in each of the described species of Edaphosaurus and show a trend for larger and more elaborate (but fewer) projecting processes over time. The possible function (or functions) of the bony tubercles on the spines remains uncertain. Romer and Price suggested that the projections on the spines of Edaphosaurus might have been embedded in tissue under the skin and might have supported food-storage or fat similar to the hump of a camel. Bennett argued that the bony projections on Edaphosaurus spines were exposed and could create air turbulence for more efficient cooling over the surface of the sail to regulate body temperature. Recent research that examined the microscopic bone structure of the tall neural spines in edaphosaurids has raised doubts about a thermoregulatory role for the sail and suggests that a display function is more plausible.
Growth and metabolism
A study comparing the microscopic bone histology of the vertebral centra of Edaphosaurus and Dimetrodon found that the plant-eating Edaphosaurus "grew distinctly more slowly" than the predator Dimetrodon, which had a higher growth rate, reflecting an "elevated metabolism". Earlier studies of Edaphosaurus limb bones had also indicated slower growth and a lower metabolism, reflecting an ectothermic (cold-blooded) animal, although the plant-eating early synapsid caseids had a lower growth rate than Edaphosaurus. Evidence of growth rates include the number of blood vessels in the bones (with more vascularization in the rapidly growing Dimetrodon) and the presence of lamellar bone in the cancellous part. In contrast to slow growth in overall body size and in most bones, the histology of the tall dorsal spines on Edaphosaurus suggests that the projecting bony tubercles developed "by sudden, rapid growth over a few seasons", unlike the incremental growth of the tubercles in the earlier edaphosaurid Ianthasaurus.
Species
Discovery and classification
Edward Drinker Cope named and described Edaphosaurus ("pavement lizard") in 1882, based on a crushed skull and a left lower jaw from the Texas Red Beds. He noted in particular the "dense body of teeth" on both the upper and lower jaws, and used the term "dental pavement" in a table in his description. The type species name pogonias means "bearded" in Greek, referring to the enlarged inward sloping chin on the lower jaw. Cope classified Edaphosaurus as a member of his Pelycosauria and created the new family Edaphosauridae. The type material did not include any of the post-cranial skeleton apart from an axis vertebra and Cope was unaware of the animal's large sail, a feature then known only for Dimetrodon.
In 1886, Cope erected the new genus Naosaurus "ship lizard" (from Greek naos "ship") for skeletal remains similar to those of the long-spined Dimetrodon, but with distinctive "transverse processes or branches, which resemble the yardarms of a ship's mast". He speculated that "the yardarms were connected by membranes with the neural spine or mast, thus serving the animal as a sail with which he navigated the waters of the Permian lakes". He recognized three species: Naosaurus claviger "club-bearer" (for the projections on its spines; now considered a synonym of Edaphosaurus pogonias); Naosaurus cruciger "cross-bearer" (for the projections on its spines; first described by Cope as Dimetrodon cruciger in 1878; now Edaphosaurus cruciger, the largest species in size); and Naosaurus microdus "small tooth" (first described as Edaphosaurus microdus in 1884). Cope noted some incomplete skull material found associated with the specimens of N. claviger and N. microdus, but thought Naosaurus was distinct from Edaphosaurus. He later decided that Naosaurus must have had a large carnivorous skull similar to Dimetrodon, although he had no direct fossil proof. In 1910, German paleontologist Otto Jaekel reported remains near Dresden in Saxony, which he called Naosaurus credneri.
In 1907, American paleontologist Ermine Cowles Case suggested that the skull of Edaphosaurus might belong with skeletons called Naosaurus, based on a specimen found in 1906 that appeared to associate elements of both. In 1913, Samuel Wendell Williston and Case described the new species Edaphosaurus novomexicanus from a fairly complete specimen unearthed in New Mexico in 1910, in which a sailbacked Naosaurus-type skeleton was found with a small Edaphosaurus-type skull. The older generic name Edaphosaurus Cope, 1882 became the valid one.
In 1940, paleontologists Alfred Sherwood Romer and Llewellyn Ivor Price named the new species Edaphosaurus boanerges ("thunderous orator") – an ironic reference to the remarkably small size of the holotype lower jaw on a composite skeleton originally mounted in the Museum of Comparative Zoology (Harvard University) with the head restored based on the larger species Edaphosaurus cruciger.
In 1979, paleontologist David Berman erected Edaphosaurus colohistion ("stunted sail") for an early species with a relatively small sail, based on fossils from West Virginia.
Reassigned species
Other proposed species of Edaphosaurus have been based on more fragmentary material that cannot be rigorously diagnosed to a genus/species level, but which may nonetheless represent edaphosaurids. The nominal species Naosaurus raymondi was assigned to Edaphosaurus by Romer and Price (1940), but Modesto and Reisz (1990) designated it a nomen vanum, and Spindler (2015) considered it probably referable to Ianthasaurus due to its age and stratigraphy. The taxon Naosaurus mirabilis Fritsch, 1895 from the Czech Republic was given its own genus Bohemiclavulus by Spindler et al. (2019).
In popular culture
The strange appearance of Edaphosaurus with its distinctive dorsal sail composed of tall spines studded with bony knobs has made it a popular subject for scientific reconstructions and paleoart in museums and in books. However, confusion over the animal's skull dating back to Cope's ideas about "Naosaurus" and over other details led to a long history of scientific and artistic errors that lasted in some cases into the 1940s. The correct scientific name Edaphosaurus (rather than "Naosaurus") also was not used consistently until the 1940s.
At the urging of paleontologist Henry Fairfield Osborn, American paleoartist Charles R. Knight consulted with Edward Drinker Cope in person in early 1897 about a set of illustrations of prehistoric reptiles, one of Cope's specialties. Shortly after, Knight reconstructed Edaphosaurus (as "Naosaurus") with a Dimetrodon skull that Cope had previously referred to that genus in error. This painting was commissioned for the American Museum of Natural History in 1897 and was reprinted for Cope's obituary in the November 1898 issue of The Century Magazine. Knight later created a more accurate revised version of the painting that turned "Naosaurus" into Dimetrodon, with a corrected head and teeth, and a sail with smooth, unbarred spines. He also turned the Dimetrodon in the original background into Edaphosaurus (still called "Naosaurus" at the time) with a different head and a sail with crossbars.
German paleontologist Otto Jaekel argued in 1905 that there was no direct scientific evidence that the tall dorsal spines on Dimetrodon and "Naosaurus" were bound in a web of skin like a sail or fin (as portrayed by Cope, Knight, and others) and proposed instead that the long bony projections served as an array of separated spines to protect the animals, which allegedly could roll up like hedgehogs. Spiny-backed reconstructions of "Naosaurus" (with a large carnivore's head) appeared in different German sources, including as a tile mosaic on the façade of the Aquarium Berlin in 1913 (destroyed in World War II and later recreated).
Nearly complete specimens of Dimetrodon and Edaphosaurus (as "Naosaurus") had not been found yet by the first decade of the 20th century when American paleontologist E.C. Case produced his major monograph on the Pelycosauria in 1907. Case argued that the apparent lack of any associated elongate and cylindrical tail bones with the known fossils meant that Dimetrodon and "Naosaurus" must have had short tails in life. (Earlier, Cope had assumed that the animals had long tails as in most reptiles, an idea seen from his sketches and his advice to Charles R. Knight in 1897.) Based on the authority of Case, museums and artists at the time restored "Naosaurus" with a short tail. New fossil finds and research by A.S. Romer in the 1930s and 1940s showed that both Dimetrodon and Edaphosaurus had long tails, a feature similar to other "pelycosaurs" and seen as primitive.
The American Museum of Natural History mounted the first full skeletal reconstruction of Edaphosaurus as "Naosaurus claviger" (a synonym of Edaphosaurus pogonias) for public display in 1907 under the scientific direction of H.F. Osborn, along with W.D. Matthew. The main part of the "Naosaurus" skeleton was a set of dorsal vertebrae with high spines (AMNH 4015) from a partial Edaphosaurus pogonias specimen found by the fossil collector Charles H. Sternberg in Hog Creek, Texas in 1896. Because of the still incomplete knowledge of Edaphosaurus at the time, the rest of the mount was a "conjectural" composite of various real fossil bones collected in different places with other parts recreated in plaster, including a skull (AMNH 4081) based on Dimetrodon (per E.D. Cope, and despite Case's already expressed doubts about such a skull for "Naosaurus") and a hypothetical short tail (per Case). As "Naosaurus" was thought to be a close relative of Dimetrodon rather than Edaphosaurus, slender limbs (AMNH 4057) probably belonging to Dimetrodon dollovianus were also mounted with this composite specimen, rather than the correct, stockier limbs now known for Edaphosaurus. The big Dimetrodon-derived skull on the museum skeleton was later replaced with one modeled on Edaphosaurus cruciger, based on more updated research. The museum eventually dismantled the entire composite restoration and by the 1950s only displayed the original set of Edaphosaurus pogonias sail vertebrae alone on the wall in Brontosaur Hall next to an accurate, fully mounted fossil skeleton of the smaller species Edaphosaurus boanerges (a nearly complete specimen (AMNH 7003) collected from Archer County, Texas, by A.S. Romer in 1939). The fossil Edaphosaurus pogonias sail spines (AMNH 4015) were remounted in the 1990s with a recreated skull (but without other skeletal parts) in a metal armature shaped in the outline of the entire animal as part of the new Hall of Primitive Mammals, which opened at the American Museum of Natural History in 1996 after major renovations.
Charles R. Knight had produced a small sculpture of a living "Naosaurus" in 1907 based on the speculative American Museum of Natural History mount. The model retained a Dimetrodon-like flesh-eater's head but differed from his earlier 1897 painted reconstruction in having a curved shape to the sail and a short tail.
The May 4, 1907 issue of Scientific American
featured a cover painting by Knight depicting a revised version of "Naosaurus" and an article (pages 368 and 370) entitled "Naosaurus: a Fossil Wonder", which described the restoration of the composite skeleton at the American Museum of Natural History and the creation of Knight's model, both under Osborn's direction.
The inaccuracy of much of Osborn's composite reconstruction of "Naosaurus" was detailed by E.C. Case in 1914 with a revised description of Edaphosaurus based on additional fossil material, including large parts of a skeleton with limb bones and a crushed skull, which Case had discovered in Archer County, Texas, in 1912 and brought to the University of Michigan. His reconstruction of Edaphosaurus cruciger, as shown in a drawing, had a much smaller head (with teeth for crushing mollusks or plants), more robust limbs, and a somewhat longer tail than Osborn's carnivorous "Naosaurus" mount. Case also confirmed that Edaphosaurus was the valid name rather than "Naosaurus". Despite his corrections, the name "Naosaurus", and even the outdated and incorrect Dimetrodon-like head, continued to appear in some popular sources.
In 1926, the Field Museum of Natural History in Chicago hired Charles R. Knight to create a series of 28 murals
(worked on from 1926 through 1930) to depict life reconstructions of prehistoric animals in the different sections of the new fossil hall of the museum for Life Over Time. One of the large murals depicted the Permian Period, with a group of five Dimetrodons, and a single Edaphosaurus, along with a group of Casea, basking in the sun surrounded by a large marsh. The Permian mural was finished in 1930. Paleontologist Elmer Riggs described the new artistic addition in the March 1931 issue of the Field Museum News and used the name "Naosaurus" for Edaphosaurus, described as "inoffensive, and given to feeding on plants". Knight's 1930 depiction of Edaphosaurus, apart from its shortened tail, was much more accurate than his earlier images of "Naosaurus" for the American Museum of Natural History, incorporating a small head and a curved profile to the sail spines.
Artist Rudolph Zallinger depicted Edaphosaurus in a more scientifically updated form (with a long tail) alongside Dimetrodon and Sphenacodon to represent the Permian period in his famous The Age of Reptiles mural (1943-1947) at the Yale Peabody Museum. The mural was based on a smaller model version of the painting in egg tempera that later appeared in The World We Live In series published in Life magazine in 1952 to 1954. The September 7, 1953 issue of Life presented The Age of Reptiles in reverse image (earliest to latest, left to right) of the mural order as a double-sided foldout page in which Edaphosaurus appeared in an Early Permian landscape with plants and animals of the period. The magazine series was edited into a popular book in 1955 that also had a foldout page for Zallinger's The Age of Reptiles artwork.
The Czech illustrator and paleoartist Zdeněk Burian created a number of vivid paintings of Edaphosaurus set in Paleozoic landscapes. (The choice to portray Edaphosaurus was based in part on edaphosaurid fossils found in native Carboniferous rocks in what is now the Czech Republic, originally identified as "Naosaurus" and now called Bohemiclavulus.) These images appeared in the series of popular general audience books on prehistoric animals that Burian produced in collaboration with Czech paleontologists Josef Augusta and Zdeněk Špinar beginning in the 1930s and on into the 1970s. Some of the books were translated into other languages, including English. Burian's painting from 1941 restored Edaphosaurus with a large carnivorous head and short tail, reflecting an outdated "Naosaurus" concept of the animal. The artwork was featured in Josef Augusta's Divy prasvěta (Wonders of the Prehistoric World), published during World War II in biweekly pamphlet form between 1941 and 1942, and then republished as a full book after the war. Burian subsequently corrected his 1941 Edaphosaurus reconstruction in a painting with the more accurate small head of a plant-eater and a long tail, the version of Edaphosaurus that appeared in later translated editions of Burian's books with Augusta such as Prehistoric Animals (1956). Another painting of Edaphosaurus by Burian appeared on the cover of the 1968 third edition of the juvenile popular science book Ztracený svět (The Lost World), also written by Augusta. The book Life Before Man (1972), written by Zdeněk Špinar, included an additional depiction of Edaphosaurus by Burian.
| Biology and health sciences | Proto-mammals | Animals |
2607409 | https://en.wikipedia.org/wiki/Laboratory%20rat | Laboratory rat | Laboratory rats or lab rats are strains of the rat subspecies Rattus norvegicus domestica (Domestic Norwegian rat) which are bred and kept for scientific research. While less commonly used for research than laboratory mice, rats have served as an important animal model for research in psychology and biomedical science, and "lab rat" is commonly used as an idiom for a test subject.
Origins of rat breeding
In 18th-century Europe, wild brown rats (Rattus norvegicus) ran rampant and this infestation fueled the industry of rat-catching. Rat-catchers would not only make money by trapping the rodents, but also by selling them for food or, more commonly, for rat-baiting.
Rat-baiting was a popular sport, which involved filling a pit with rats and timing how long it took for a terrier to kill them all. Over time, breeding the rats for these contests may have produced variations in color, notably the albino and hooded varieties. The first time one of these albino mutants was brought into a laboratory for a study was in 1828 for an experiment on fasting. Over the next 30 years, rats were used for several more experiments and eventually the laboratory rat became the first animal domesticated for purely scientific reasons.
In Japan, there was a widespread practice of keeping rats as a domesticated pet during the Edo period and in the 18th century guidebooks on keeping domestic rats were published by Youso Tamanokakehashi (1775) and Chingan Sodategusa (1787). Genetic analysis of 117 albino rat strains collected from all parts of the world carried out by a team led by Takashi Kuramoto at Kyoto University in 2012 showed that the albinos descended from hooded rats and all the albinos descended from a single ancestor. As there is evidence that the hooded rat was known as the "Japanese rat" in the early 20th century, Kuramoto concluded that one or more Japanese hooded rats might have been brought to Europe or the Americas and an albino rat that emerged as a product of the breeding of these hooded rats was the common ancestor of all the albino laboratory rats in use today.
Use in research
The rat found early use in laboratory research in five areas: W. S. Small suggested that the rate of learning could be measured by rats in a maze; a suggestion employed by John B. Watson for his Ph.D. dissertation in 1903. The first rat colony in America used for nutrition research was started in January 1908 by Elmer McCollum and then, nutritive requirements of rats were used by Thomas Burr Osborne and Lafayette Mendel to determine the details of protein nutrition. The reproductive function of rats was studied at the Institute for Experimental Biology at the University of California, Berkeley by Herbert McLean Evans and Joseph A. Long. The genetics of rats was studied by William Ernest Castle at the Bussey Institute of Harvard University until it closed in 1994. Rats have long been used in cancer research; for instance at the Crocker Institute for Cancer Research.
The historical importance of this species to scientific research is reflected by the amount of literature on it: roughly 50% more than that on laboratory mice. Laboratory rats are frequently subject to dissection or microdialysis to study internal effects on organs and the brain, such as for cancer or pharmacological research. Laboratory rats not sacrificed may be euthanized or, in some cases, become pets.
Domestic rats differ from wild rats (various spp. of Rodentia) in many ways: they are calmer and significantly less likely to bite, they can tolerate greater crowding, they breed earlier and produce more offspring, and their brains, livers, kidneys, adrenal glands, and hearts are smaller.
Scientists have bred many strains or "lines" of rats specifically for experimentation. Most are derived from the albino Wistar rat, which is still widely used. Other common strains are the Sprague Dawley, Fischer 344, Holtzman albino strains, Long–Evans, and Lister black hooded rats. Inbred strains are also available, but are not as commonly used as inbred mice.
Much of the genome of Rattus norvegicus has been sequenced. In October 2003, researchers succeeded in cloning two laboratory rats by nuclear transfer. This was the first in a series of developments that have begun to make rats tractable as genetic research subjects, although they still lag behind mice, which lend themselves better to the embryonic stem cell techniques typically used for genetic manipulation. Many investigators who wish to trace observations on behavior and physiology to underlying genes regard aspects of these in rats as more relevant to humans and easier to observe than in mice, giving impetus to the development of genetic research techniques applicable to rats.
A 1972 study compared neoplasms in Sprague Dawleys from six different commercial suppliers and found highly significant differences in the incidences of endocrine and mammary tumors. There were even significant variations in the incidences of adrenal medulla tumors among rats from the same source raised in different laboratories. All but one of the testicular tumors occurred in the rats from a single supplier. The researchers found that the incidence of tumors in Sprague Dawleys from different suppliers varied as much from each other as from the other strains of rats. The authors of the study "stressed the need for extreme caution in evaluation of carcinogenicity studies conducted at different laboratories and/or on rats from different sources."
During food rationing due to World War II, British biologists had eaten laboratory rats, creamed.
Scientists have also spent time studying the thermoregulation of the rat's tail in research. The rat's tail works as a variable heat exchanger. The tail's blood flow allows for thermoregulation to take place because it is under control of sympathetic vasoconstrictor nerves. Vasodilation occurs when the tail temperature increases, causing heat loss. Vasoconstriction occurs when the tail temperature decreases allowing heat to be conserved. Thermoregulation in the rat tail has been used to study metabolism.
Stocks and strains
A "strain", in reference to rodents, is a group in which all members are, as nearly as possible, genetically identical. In rats, this is accomplished through inbreeding. By having this kind of population, it is possible to conduct experiments on the roles of genes, or conduct experiments that exclude variations in genetics as a factor. By contrast, "outbred" populations are used when identical genotypes are unnecessary or a population with genetic variation is required, and these rats are usually referred to as "stocks" rather than "strains".
Wistar rat
The Wistar rat is an outbred albino rat. This breed was developed at the Wistar Institute in 1906 for use in biological and medical research, and is notably the first rat developed to serve as a model organism at a time when laboratories primarily used the house mouse (Mus musculus). More than half of all laboratory rat strains are descended from the original colony established by physiologist Henry Herbert Donaldson, scientific administrator Milton J. Greenman, and genetic researcher/embryologist Helen Dean King.
The Wistar rat is currently one of the most popular rats used for laboratory research. It is characterized by its wide head, long ears, and a tail length that is always less than its body length. The Sprague Dawley and Long–Evans were developed from Wistars. Wistars are more active than others like Sprague Dawleys. The spontaneously hypertensive rat and the Lewis are other well-known stocks developed from Wistars.
Long–Evans rat
The Long–Evans rat is an outbred rat developed by Long and Evans in 1915 by crossbreeding several Wistar females with a wild gray male. Long-Evans rats are white with a black hood, or occasionally white with a brown hood. They are utilized as a multipurpose model organism, frequently in behavioral research, especially in alcohol research. Long-Evans consume alcohol in a much higher rate compared to other strains, thus require less time for these behavioral studies.
Sprague Dawley rat
The Sprague Dawley is an outbred, multipurpose breed of albino rat used extensively in medical and nutritional research. Its main advantage is its calmness and ease of handling. This breed of rat was first produced by the Sprague Dawley farms (later to become the Sprague Dawley Animal Company) in Madison, Wisconsin, in 1925. The name was originally hyphenated, although the brand styling today (Sprague Dawley, the trademark used by Inotiv) is not. The average litter size of the Sprague Dawley rat is 11.0.
These rats typically have a longer tail in proportion to their body length than Wistars. They were used in the Séralini affair, where the herbicide RoundUp was claimed to increase the occurrence of tumor in these rats. However, since these rats are known to grow tumors at a high (and very variable) rate, the study was considered flawed in design and its findings unsubstantiated.
Biobreeding rat
The biobreeding rat (a.k.a. the biobreeding diabetes-prone rat or BBDP rat) is an inbred strain that spontaneously develops autoimmune type 1 diabetes. Like NOD mice, biobreeding rats are used as an animal model for Type 1 diabetes. The strain re-capitulates many of the features of human type 1 diabetes and has contributed greatly to the research of T1DM pathogenesis.
Brattleboro rat
The Brattleboro rat is a strain that was developed by Henry A. Schroeder and technician Tim Vinton in West Brattleboro, Vermont, beginning in 1961, for Dartmouth Medical School. It has a naturally occurring genetic mutation that makes specimens unable to produce the hormone vasopressin, which helps control kidney function. The rats were being raised for laboratory use by Henry Schroeder and technician Tim Vinton, who noticed that the litter of 17 drank and urinated excessively.
Hairless rat
Hairless laboratory rats provide researchers with valuable data regarding compromised immune systems and genetic kidney diseases. It is estimated that there are over 25 genes that cause recessive hairlessness in laboratory rats. The more common ones are denoted as rnu (Rowett nude), fz (fuzzy), and shn (shorn).
Rowett nude rats, first identified in 1953 in Scotland, have no thymus. The lack of this organ severely compromises their immune system, with infections of the respiratory tract and eyes increasing the most dramatically.
Fuzzy rats were identified in 1976 in a Pennsylvania lab. The leading cause of death among fz/fz rats is ultimately a progressive kidney failure that begins around the age of 1 year.
Shorn rats were bred from Sprague Dawley rats in Connecticut in 1998. They also suffer from severe kidney problems.
Lewis rat
The Lewis rat was developed by Margaret Lewis from Wistar stock in the early 1950s. Characteristics include albino coloring, docile behavior, and low fertility.
The Lewis rat suffers from several spontaneous pathologies: first, they can suffer from high incidences of neoplasms, with the rat's lifespan mainly determined by this. The most common are adenomas of the pituitary and adenomas/adenocarcinomas of the adrenal cortex in both sexes, mammary gland tumors and endometrial carcinomas in females, and C-cell adenomas/adenocarcinomas of the thyroid gland and tumors of the hematopoietic system in males. Second, Lewis rats are prone to develop a spontaneous transplantable lymphatic leukaemia. Lastly, when in advanced age, they sometimes develop spontaneous glomerular sclerosis.
Research applications include transplantation research, induced arthritis and inflammation, experimental allergic encephalitis, and STZ-induced diabetes.
Royal College of Surgeons rat
The Royal College of Surgeons rat (or RCS rat) is the first known animal with inherited retinal degeneration. Although the genetic defect was not known for many years, it was identified in the year 2000 as a mutation in the gene MERTK. This mutation results in defective retinal pigment epithelium phagocytosis of photoreceptor outer segments.
Shaking rat Kawasaki
The shaking rat Kawasaki (SRK) is an autosomal recessive mutant that has a short deletion in the RELN (reelin) gene. This results in the lowered expression of reelin protein, essential for proper cortex lamination and cerebellum development. Its phenotype is similar to the widely researched reeler mouse. Shaking rat Kawasaki was first described in 1988. This and the Lewis rat are well-known stocks developed from Wistar rats.
Zucker rat
The Zucker rat was bred to be a genetic model for research on obesity and hypertension. They are named after Lois M. Zucker and Theodore F. Zucker, pioneer researchers in the study of the genetics of obesity. There are two types of Zucker rat: a lean Zucker rat, denoted as the dominant trait (Fa/Fa) or (Fa/fa); and the characteristically obese (or fatty) Zucker rat or Zucker diabetic fatty rat (ZDF rat), which is actually a recessive trait (fa/fa) of the leptin receptor, capable of weighing up to — more than twice the average weight.
Obese Zucker rats have high levels of lipids and cholesterol in their bloodstream, are resistant to insulin without being hyperglycemic, and gain weight from an increase in both the size and number of fat cells. Obesity in Zucker rats is primarily linked to their hyperphagic nature and excessive hunger; however, food intake does not fully explain the hyperlipidemia or overall body composition.
Knockout rats
A knockout rat (also spelled knock out or knock-out) is a genetically engineered rat with a single gene turned off through a targeted mutation. Knockout rats can mimic human diseases, and are important tools for studying gene function and for drug discovery and development. The production of knockout rats became technically feasible in 2008, through work financed by $120 million in funding from the National Institutes of Health (NIH) via the Rat Genome Sequencing Project Consortium, and work accomplished by the members of the Knock Out Rat Consortium (KORC). Knockout rat disease models for Parkinson's disease, Alzheimer's disease, hypertension, and diabetes, using zinc-finger nuclease technology, are being commercialized by SAGE Labs.
| Biology and health sciences | Rodents | Animals |
2608645 | https://en.wikipedia.org/wiki/Polybolos | Polybolos | The polybolos (the name means "multi-thrower" in Greek) was an ancient Greek repeating ballista, reputedly invented by Dionysius of Alexandria (a 3rd-century BC Greek engineer at the Rhodes arsenal,) and used in antiquity. The polybolos was not a crossbow since it used a torsion mechanism, drawing its power from twisted sinew-bundles. However the earlier and similar oxybeles employed a tension crosbow mechanism, before it was abandoned in favor of torsion.
Philo of Byzantium ( 280 BC – 220 BC) encountered and described a weapon similar to the polybolos, a catapult that could fire again and again without a need for manual reloading. Philo left a detailed description of the gears that powered its chain drive (the oldest known application of such a mechanism) and that placed bolt after bolt into its firing slot.
Design
The polybolos would have differed from an ordinary ballista in that it had a wooden hopper magazine, capable of holding several dozen bolts, that was positioned over the mensa (the cradle that holds the bolt prior to firing). The mechanism is unique in that it is driven by a flat-link chain connected to a windlass. The mensa itself was a sliding plank (similar to that on the gastraphetes) containing the claw latches used to pull back the drawstring and was attached to the chain link. When loading a new bolt and spanning the drawstring, the windlass is rotated counterclockwise by an operator standing on the left side of the weapon; this drives the mensa forward towards the bow string. At the very front, a metal lug triggers the latching claws into catching the drawstring.
Once the string is held firm by the trigger mechanism, the windlass is then rotated clockwise; pulling the mensa back and drawing the bow string with it. At the same time, a round wooden pole in the bottom of the magazine is rotated via a spiral groove being driven by a rivet attached to the sliding mensa; dropping a single bolt from a carved notch in the rotating pole. With the drawstring pulled back and a bolt loaded on the mensa, the polybolos is ready to be fired. As the windlass is rotated further back to the very back end, the claws on the mensa meets another lug like the one that pushed the claws into catching the string. This one causes the claws to disengage the drawstring and automatically fires the loaded bolt. Upon the bolt being fired, the process is repeated. The repetition provides the weapon's name, in Greek , "throwing many missiles", from (), "multiple, many" and () "thrower", in turn from (), "to throw, to hurl", literally a multithrower.
Popular culture
In 2010, a MythBusters episode was dedicated to building and testing a replica, and concluded that its existence as a historical weapon was plausible. However, the machine MythBusters built was prone to breakdowns that had to be fixed multiple times.
| Technology | Artillery | null |
2609001 | https://en.wikipedia.org/wiki/Hungarian%20algorithm | Hungarian algorithm | The Hungarian method is a combinatorial optimization algorithm that solves the assignment problem in polynomial time and which anticipated later primal–dual methods. It was developed and published in 1955 by Harold Kuhn, who gave it the name "Hungarian method" because the algorithm was largely based on the earlier works of two Hungarian mathematicians, Dénes Kőnig and Jenő Egerváry. However, in 2006 it was discovered that Carl Gustav Jacobi had solved the assignment problem in the 19th century, and the solution had been published posthumously in 1890 in Latin.
James Munkres reviewed the algorithm in 1957 and observed that it is (strongly) polynomial. Since then the algorithm has been known also as the Kuhn–Munkres algorithm or Munkres assignment algorithm. The time complexity of the original algorithm was , however Edmonds and Karp, and independently Tomizawa, noticed that it can be modified to achieve an running time. Ford and Fulkerson extended the method to general maximum flow problems in form of the Ford–Fulkerson algorithm.
The problem
Example
In this simple example, there are three workers: Alice, Bob and Carol. One of them has to clean the bathroom, another sweep the floors and the third washes the windows, but they each demand different pay for the various tasks. The problem is to find the lowest-cost way to assign the jobs. The problem can be represented in a matrix of the costs of the workers doing the jobs. For example:
{| class="wikitable" style="text-align:center;"
!
! Cleanbathroom
! Sweepfloors
! Wash windows
|-
! Alice
| $8
| $4
| $7
|-
! Bob
| $5
| $2
| $3
|-
! Carol
| $9
| $4
| $8
|}
The Hungarian method, when applied to the above table, would give the minimum cost: this is $15, achieved by having Alice clean the bathroom, Carol sweep the floors, and Bob wash the windows. This can be confirmed using brute force:
{| class="wikitable" style="text-align:center;"
!
! Alice
! Bob
! Carol
|-
! Alice
| —
| $17
| $16
|-
! Bob
| $18
| —
| $18
|-
! Carol
| $15
| $16
| —
|} (the unassigned person washes the windows)
Matrix formulation
In the matrix formulation, we are given an n×n matrix, where the element in the i-th row and j-th column represents the cost of assigning the j-th job to the i-th worker. We have to find an assignment of the jobs to the workers, such that each job is assigned to one worker and each worker is assigned one job, such that the total cost of assignment is minimum.
This can be expressed as permuting the rows of a cost matrix C to minimize the trace of a matrix,
where P is a permutation matrix. (Equivalently, the columns can be permuted using CP.)
If the goal is to find the assignment that yields the maximum cost, the problem can be solved by negating the cost matrix C.
Bipartite graph formulation
The algorithm can equivalently be described by formulating the problem using a bipartite graph. We have a complete bipartite graph with worker vertices () and job vertices (), and the edges () each have a cost . We want to find a perfect matching with a minimum total cost.
The algorithm in terms of bipartite graphs
Let us call a function a potential if for each .
The value of potential is the sum of the potential over all vertices:
.
The cost of each perfect matching is at least the value of each potential: the total cost of the matching is the sum of costs of all edges; the cost of each edge is at least the sum of potentials of its endpoints; since the matching is perfect, each vertex is an endpoint of exactly one edge; hence the total cost is at least the total potential.
The Hungarian method finds a perfect matching and a potential such that the matching cost equals the potential value. This proves that both of them are optimal. In fact, the Hungarian method finds a perfect matching of tight edges: an edge is called tight for a potential if . Let us denote the subgraph of tight edges by . The cost of a perfect matching in (if there is one) equals the value of .
During the algorithm we maintain a potential and an orientation of (denoted by ) which has the property that the edges oriented from to form a matching . Initially, is 0 everywhere, and all edges are oriented from to (so is empty). In each step, either we modify so that its value increases, or modify the orientation to obtain a matching with more edges. We maintain the invariant that all the edges of are tight. We are done if is a perfect matching.
In a general step, let and be the vertices not covered by (so consists of the vertices in with no incoming edge and consists of the vertices in with no outgoing edge). Let be the set of vertices reachable in from by a directed path. This can be computed by breadth-first search.
If is nonempty, then reverse the orientation of all edges along a directed path in from to . Thus the size of the corresponding matching increases by 1.
If is empty, then let
is well defined because at least one such edge must exist whenever the matching is not yet of maximum possible size (see the following section); it is positive because there are no tight edges between and . Increase by on the vertices of and decrease by on the vertices of . The resulting is still a potential, and although the graph changes, it still contains (see the next subsections). We orient the new edges from to . By the definition of the set of vertices reachable from increases (note that the number of tight edges does not necessarily increase). If the vertex added to is unmatched (that is, it is also in ), then at the next iteration the graph will have an augmenting path.
We repeat these steps until is a perfect matching, in which case it gives a minimum cost assignment. The running time of this version of the method is : is augmented times, and in a phase where is unchanged, there are at most potential changes (since increases every time). The time sufficient for a potential change is .
Proof that the algorithm makes progress
We must show that as long as the matching is not of maximum possible size, the algorithm is always able to make progress — that is, to either increase the number of matched edges, or tighten at least one edge. It suffices to show that at least one of the following holds at every step:
is of maximum possible size.
contains an augmenting path.
contains a loose-tailed path: a path from some vertex in to a vertex in that consists of any number (possibly zero) of tight edges followed by a single loose edge. The trailing loose edge of a loose-tailed path is thus from , guaranteeing that is well defined.
If is of maximum possible size, we are of course finished. Otherwise, by Berge's lemma, there must exist an augmenting path with respect to in the underlying graph . However, this path may not exist in : Although every even-numbered edge in is tight by the definition of , odd-numbered edges may be loose and thus absent from . One endpoint of is in , the other in ; w.l.o.g., suppose it begins in . If every edge on is tight, then it remains an augmenting path in and we are done. Otherwise, let be the first loose edge on . If then we have found a loose-tailed path and we are done. Otherwise, is reachable from some other path of tight edges from a vertex in . Let be the subpath of beginning at and continuing to the end, and let be the path formed by traveling along until a vertex on is reached, and then continuing to the end of . Observe that is an augmenting path in with at least one fewer loose edge than . can be replaced with and this reasoning process iterated (formally, using induction on the number of loose edges) until either an augmenting path in or a loose-tailed path in is found.
Proof that adjusting the potential y leaves M unchanged
To show that every edge in remains after adjusting , it suffices to show that for an arbitrary edge in , either both of its endpoints, or neither of them, are in . To this end let be an edge in from to . It is easy to see that if is in then must be too, since every edge in is tight. Now suppose, toward contradiction, that but . itself cannot be in because it is the endpoint of a matched edge, so there must be some directed path of tight edges from a vertex in to . This path must avoid , since that is by assumption not in , so the vertex immediately preceding in this path is some other vertex . is a tight edge from to and is thus in . But then contains two edges that share the vertex , contradicting the fact that is a matching. Thus every edge in has either both endpoints or neither endpoint in .
Proof that remains a potential
To show that remains a potential after being adjusted, it suffices to show that no edge has its total potential increased beyond its cost. This is already established for edges in by the preceding paragraph, so consider an arbitrary edge from to . If is increased by , then either , in which case is decreased by , leaving the total potential of the edge unchanged, or , in which case the definition of guarantees that . Thus remains a potential.
The algorithm in O(n3) time
Suppose there are jobs and workers (). We describe how to compute for each prefix of jobs the minimum total cost to assign each of these jobs to distinct workers. Specifically, we add the th job and update the total cost in time , yielding an overall time complexity of . Note that this is better than when the number of jobs is small relative to the number of workers.
Adding the j-th job in O(jW) time
We use the same notation as the previous section, though we modify their definitions as necessary. Let denote the set of the first jobs and denote the set of all workers.
Before the th step of the algorithm, we assume that we have a matching on that matches all jobs in and potentials satisfying the following condition: the matching is tight with respect to the potentials, and the potentials of all unmatched workers are zero, and the potentials of all matched workers are non-positive. Note that such potentials certify the optimality of the matching.
During the th step, we add the th job to to form and initialize . At all times, every vertex in will be reachable from the th job in . While does not contain a worker that has not been assigned a job, let
and denote any at which the minimum is attained. After adjusting the potentials in the way described in the previous section, there is now a tight edge from to .
If is unmatched, then we have an augmenting path in the subgraph of tight edges from to . After toggling the matching along this path, we have now matched the first jobs, and this procedure terminates.
Otherwise, we add and the job matched with it to .
Adjusting potentials takes time. Recomputing and after changing the potentials and also can be done in time. Case 1 can occur at most times before case 2 occurs and the procedure terminates, yielding the overall time complexity of .
Implementation in C++
For convenience of implementation, the code below adds an additional worker such that stores the negation of the sum of all computed so far. After the th job is added and the matching updated, the cost of the current matching equals the sum of all computed so far, or .
This code is adapted from e-maxx :: algo.
/**
* Solution to https://open.kattis.com/problems/cordonbleu using Hungarian
* algorithm.
*/
#include <cassert>
#include <iostream>
#include <limits>
#include <vector>
using namespace std;
/**
* Sets a = min(a, b)
* @return true if b < a
*/
template <class T> bool ckmin(T &a, const T &b) { return b < a ? a = b, 1 : 0; }
/**
* Given J jobs and W workers (J <= W), computes the minimum cost to assign each
* prefix of jobs to distinct workers.
*
* @tparam T a type large enough to represent integers on the order of J *
* max(|C|)
* @param C a matrix of dimensions JxW such that C[j][w] = cost to assign j-th
* job to w-th worker (possibly negative)
*
* @return a vector of length J, with the j-th entry equaling the minimum cost
* to assign the first (j+1) jobs to distinct workers
*/
template <class T> vector<T> hungarian(const vector<vector<T>> &C) {
const int J = (int)size(C), W = (int)size(C[0]);
assert(J <= W);
// job[w] = job assigned to w-th worker, or -1 if no job assigned
// note: a W-th worker was added for convenience
vector<int> job(W + 1, -1);
vector<T> ys(J), yt(W + 1); // potentials
// -yt[W] will equal the sum of all deltas
vector<T> answers;
const T inf = numeric_limits<T>::max();
for (int j_cur = 0; j_cur < J; ++j_cur) { // assign j_cur-th job
int w_cur = W;
job[w_cur] = j_cur;
// min reduced cost over edges from Z to worker w
vector<T> min_to(W + 1, inf);
vector<int> prv(W + 1, -1); // previous worker on alternating path
vector<bool> in_Z(W + 1); // whether worker is in Z
while (job[w_cur] != -1) { // runs at most j_cur + 1 times
in_Z[w_cur] = true;
const int j = job[w_cur];
T delta = inf;
int w_next;
for (int w = 0; w < W; ++w) {
if (!in_Z[w]) {
if (ckmin(min_to[w], C[j][w] - ys[j] - yt[w]))
prv[w] = w_cur;
if (ckmin(delta, min_to[w])) w_next = w;
}
}
// delta will always be nonnegative,
// except possibly during the first time this loop runs
// if any entries of C[j_cur] are negative
for (int w = 0; w <= W; ++w) {
if (in_Z[w]) ys[job[w]] += delta, yt[w] -= delta;
else min_to[w] -= delta;
}
w_cur = w_next;
}
// update assignments along alternating path
for (int w; w_cur != W; w_cur = w) job[w_cur] = job[w = prv[w_cur]];
answers.push_back(-yt[W]);
}
return answers;
}
/**
* Sanity check: https://en.wikipedia.org/wiki/Hungarian_algorithm#Example
* First job (5):
* clean bathroom: Bob -> 5
* First + second jobs (9):
* clean bathroom: Bob -> 5
* sweep floors: Alice -> 4
* First + second + third jobs (15):
* clean bathroom: Alice -> 8
* sweep floors: Carol -> 4
* wash windows: Bob -> 3
*/
void sanity_check_hungarian() {
vector<vector<int>> costs{{8, 5, 9}, {4, 2, 4}, {7, 3, 8}};
assert((hungarian(costs) == vector<int>{5, 9, 15}));
cerr << "Sanity check passed.\n";
}
// solves https://open.kattis.com/problems/cordonbleu
void cordon_bleu() {
int N, M;
cin >> N >> M;
vector<pair<int, int>> B(N), C(M);
vector<pair<int, int>> bottles(N), couriers(M);
for (auto &b : bottles) cin >> b.first >> b.second;
for (auto &c : couriers) cin >> c.first >> c.second;
pair<int, int> rest;
cin >> rest.first >> rest.second;
vector<vector<int>> costs(N, vector<int>(N + M - 1));
auto dist = [&](pair<int, int> x, pair<int, int> y) {
return abs(x.first - y.first) + abs(x.second - y.second);
};
for (int b = 0; b < N; ++b) {
for (int c = 0; c < M; ++c) { // courier -> bottle -> restaurant
costs[b][c] =
dist(couriers[c], bottles[b]) + dist(bottles[b], rest);
}
for (int _ = 0; _ < N - 1; ++_) { // restaurant -> bottle -> restaurant
costs[b][_ + M] = 2 * dist(bottles[b], rest);
}
}
cout << hungarian(costs).back() << "\n";
}
int main() {
sanity_check_hungarian();
cordon_bleu();
}
Connection to successive shortest paths
The Hungarian algorithm can be seen to be equivalent to the successive shortest path algorithm for minimum cost flow, where the reweighting technique from Johnson's algorithm is used to find the shortest paths.
The implementation from the previous section is rewritten below in such a way as to emphasize this
connection; it can be checked that the potentials for workers are equal to the potentials from the previous solution up to a constant offset. When the graph is sparse (there are only allowed job, worker pairs), it is possible
to optimize this algorithm to run in time by using a Fibonacci heap to determine
instead of iterating over all workers to find the one with minimum distance (alluded to here).
template <class T> vector<T> hungarian(const vector<vector<T>> &C) {
const int J = (int)size(C), W = (int)size(C[0]);
assert(J <= W);
// job[w] = job assigned to w-th worker, or -1 if no job assigned
// note: a W-th worker was added for convenience
vector<int> job(W + 1, -1);
vector<T> h(W); // Johnson potentials
vector<T> answers;
T ans_cur = 0;
const T inf = numeric_limits<T>::max();
// assign j_cur-th job using Dijkstra with potentials
for (int j_cur = 0; j_cur < J; ++j_cur) {
int w_cur = W; // unvisited worker with minimum distance
job[w_cur] = j_cur;
vector<T> dist(W + 1, inf); // Johnson-reduced distances
dist[W] = 0;
vector<bool> vis(W + 1); // whether visited yet
vector<int> prv(W + 1, -1); // previous worker on shortest path
while (job[w_cur] != -1) { // Dijkstra step: pop min worker from heap
T min_dist = inf;
vis[w_cur] = true;
int w_next = -1; // next unvisited worker with minimum distance
// consider extending shortest path by w_cur -> job[w_cur] -> w
for (int w = 0; w < W; ++w) {
if (!vis[w]) {
// sum of reduced edge weights w_cur -> job[w_cur] -> w
T edge = C[job[w_cur]][w] - h[w];
if (w_cur != W) {
edge -= C[job[w_cur]][w_cur] - h[w_cur];
assert(edge >= 0); // consequence of Johnson potentials
}
if (ckmin(dist[w], dist[w_cur] + edge)) prv[w] = w_cur;
if (ckmin(min_dist, dist[w])) w_next = w;
}
}
w_cur = w_next;
}
for (int w = 0; w < W; ++w) { // update potentials
ckmin(dist[w], dist[w_cur]);
h[w] += dist[w];
}
ans_cur += h[w_cur];
for (int w; w_cur != W; w_cur = w) job[w_cur] = job[w = prv[w_cur]];
answers.push_back(ans_cur);
}
return answers;
}
Matrix interpretation
This variant of the algorithm follows the formulation given by Flood, and later described more explicitly by Munkres, who proved it runs in time. Instead of keeping track of the potentials of the vertices, the algorithm operates only on a matrix:
where is the original cost matrix and are the potentials from the graph interpretation. Changing the potentials corresponds to adding or subtracting from rows or columns of this matrix. The algorithm starts with . As such, it can be viewed as taking the original cost matrix and modifying it.
Given workers and tasks, the problem is written in the form of an × cost matrix
where a, b, c and d are workers who have to perform tasks 1, 2, 3 and 4. a1, a2, a3, and a4 denote the penalties incurred when worker "a" does task 1, 2, 3, and 4 respectively.
The problem is equivalent to assigning each worker a unique task such that the total penalty is minimized. Note that each task can only be worked on by one worker.
Step 1
For each row, its minimum element is subtracted from every element in that row. This causes all elements to have nonnegative values. Therefore, an assignment with a total penalty of 0 is by definition a minimum assignment.
This also leads to at least one zero in each row. As such, a naive greedy algorithm can attempt to assign all workers a task with a penalty of zero. This is illustrated below.
The zeros above would be the assigned tasks.
Worst-case there are ! combinations to try, since multiple zeroes can appear in a row if multiple elements are the minimum. So at some point this naive algorithm should be short circuited.
Step 2
Sometimes it may turn out that the matrix at this stage cannot be used for assigning, as is the case for the matrix below.
To overcome this, we repeat the above procedure for all columns (i.e. the minimum element in each column is subtracted from all the elements in that column) and then check if an assignment with penalty 0 is possible.
In most situations this will give the result, but if it is still not possible then we need to keep going.
Step 3
All zeros in the matrix must be covered by marking as few rows and/or columns as possible. Steps 3 and 4 form one way to accomplish this.
For each row, try to assign an arbitrary zero. Assigned tasks are represented by starring a zero. Note that assignments can't be in the same row or column.
We assign the first zero of Row 1. The second zero of Row 1 can't be assigned.
We assign the first zero of Row 2. The second zero of Row 2 can't be assigned.
Zeros on Row 3 and Row 4 can't be assigned, because they are on the same column as the zero assigned on Row 1.
We could end with another assignment if we choose another ordering of the rows and columns.
Step 4
Cover all columns containing a (starred) zero.
{|class="wikitable" style="text-align:center"
|- style="background: white"
|× ||× || || ||
|-
|style="background:lightgrey"| 0*||style="background:lightgrey"|a2|| 0 ||a4
|style="background: white"|
|-
|style="background:lightgrey"|b1||style="background:lightgrey"| 0*||b3|| 0
|style="background: white"|
|-
|style="background:lightgrey"| 0 ||style="background:lightgrey"|c2||c3||c4
|style="background: white"|
|-
|style="background:lightgrey"| 0 ||style="background:lightgrey"|d2||d3||d4
|}
Find a non-covered zero and prime it (mark it with a prime symbol). If no such zero can be found, meaning all zeroes are covered, skip to step 5.
If the zero is on the same row as a starred zero, cover the corresponding row, and uncover the column of the starred zero.
Then, GOTO "Find a non-covered zero and prime it."
Here, the second zero of Row 1 is uncovered. Because there is another zero starred on Row 1, we cover Row 1 and uncover Column 1.
Then, the second zero of Row 2 is uncovered. We cover Row 2 and uncover Column 2.
{|class="wikitable" style="text-align:center"
|- style="background: white"
|| ||× || || ||
|- style="background:lightgrey"
| 0*||a2|| 0'||a4
|style="background: white"|×
|-
|b1||style="background:lightgrey"| 0*||b3|| 0
|style="background: white"|
|-
|| 0 ||style="background:lightgrey"|c2||c3||c4
|style="background: white"|
|-
|| 0 ||style="background:lightgrey"|d2||d3||d4
|style="background: white"|
|}
{|class="wikitable" style="text-align:center"
|- style="background: white"
|| || || || ||
|- style="background:lightgrey"
| 0*||a2|| 0'||a4
|style="background: white"|×
|- style="background:lightgrey"
|b1|| 0*||b3|| 0'
|style="background: white"|×
|-
|| 0 ||c2||c3||c4
|style="background: white"|
|-
|| 0 ||d2||d3||d4
|style="background: white"|
|}
Else the non-covered zero has no assigned zero on its row. We make a path starting from the zero by performing the following steps:
Substep 1: Find a starred zero on the corresponding column. If there is one, go to Substep 2, else, stop.
Substep 2: Find a primed zero on the corresponding row (there should always be one). Go to Substep 1.
The zero on Row 3 is uncovered. We add to the path the first zero of Row 1, then the second zero of Row 1, then we are done.
{|class="wikitable" style="text-align:center"
|- style="background: white"
|| || || || ||
|- style="background:lightgrey"
| 0* ||a2|| 0' ||a4
|style="background: white"|×
|- style="background:lightgrey"
|b1|| 0*||b3|| 0'
|style="background: white"|×
|-
|| 0' ||c2||c3||c4
|style="background: white"|
|-
|| 0 ||d2||d3||d4
|style="background: white"|
|}
(Else branch continued) For all zeros encountered during the path, star primed zeros and unstar starred zeros.
As the path begins and ends by a primed zero when swapping starred zeros, we have assigned one more zero.
{|class="wikitable" style="text-align:center"
|-
| 0 ||a2|| 0*||a4
|-
|b1|| 0*||b3|| 0
|-
|| 0*||c2||c3||c4
|-
|| 0 ||d2||d3||d4
|}
(Else branch continued) Unprime all primed zeroes and uncover all lines.
Repeat the previous steps (continue looping until the above "skip to step 5" is reached).
We cover columns 1, 2 and 3. The second zero on Row 2 is uncovered, so we cover Row 2 and uncover Column 2:
{|class="wikitable" style="text-align:center"
|- style="background: white"
||× || ||× || ||
|-
|style="background:lightgrey"| 0 ||a2||style="background:lightgrey"| 0*||a4
|style="background: white"|
|- style="background:lightgrey"
|b1|| 0*||b3|| 0'
|style="background: white"|×
|-
|style="background:lightgrey"| 0*||c2||style="background:lightgrey"|c3||c4
|style="background: white"|
|-
|style="background:lightgrey"| 0 ||d2||style="background:lightgrey"|d3||d4
|style="background: white"|
|}
All zeros are now covered with a minimal number of rows and columns.
The aforementioned detailed description is just one way to draw the minimum number of lines to cover all the 0s. Other methods work as well.
Step 5
If the number of starred zeros is (or in the general case , where is the number of people and is the number of jobs), the algorithm terminates. See the Result subsection below on how to interpret the results.
Otherwise, find the lowest uncovered value. Subtract this from every unmarked element and add it to every element covered by two lines. Go back to step 4.
This is equivalent to subtracting a number from all rows which are not covered and adding the same number to all columns which are covered. These operations do not change optimal assignments.
Result
If following this specific version of the algorithm, the starred zeros form the minimum assignment.
From Kőnig's theorem, the minimum number of lines (minimum vertex cover) will be (the size of maximum matching). Thus, when lines are required, minimum cost assignment can be found by looking at only zeroes in the matrix.
Bibliography
R.E. Burkard, M. Dell'Amico, S. Martello: Assignment Problems (Revised reprint). SIAM, Philadelphia (PA.) 2012.
M. Fischetti, "Lezioni di Ricerca Operativa", Edizioni Libreria Progetto Padova, Italia, 1995.
R. Ahuja, T. Magnanti, J. Orlin, "Network Flows", Prentice Hall, 1993.
S. Martello, "Jeno Egerváry: from the origins of the Hungarian algorithm to satellite communication". Central European Journal of Operational Research 18, 47–58, 2010
| Mathematics | Optimization | null |
2611033 | https://en.wikipedia.org/wiki/Mangrove%20forest | Mangrove forest | Mangrove forests, also called mangrove swamps, mangrove thickets or mangals, are productive wetlands that occur in coastal intertidal zones. Mangrove forests grow mainly at tropical and subtropical latitudes because mangrove trees cannot withstand freezing temperatures. There are about 80 different species of mangroves, all of which grow in areas with low-oxygen soil, where slow-moving waters allow fine sediments to accumulate.
Many mangrove forests can be recognised by their dense tangle of prop roots that make the trees appear to be standing on stilts above the water. This tangle of roots allows the trees to handle the daily rise and fall of tides, as most mangroves get flooded at least twice per day. The roots slow the movement of tidal waters, causing sediments to settle out of the water and build up the muddy bottom. Mangrove forests stabilise the coastline, reducing erosion from storm surges, currents, waves, and tides. The intricate root system of mangroves also makes these forests attractive to fish and other organisms seeking food and shelter from predators.
Mangrove forests live at the interface between the land, the ocean, and the atmosphere, and are centres for the flow of energy and matter between these systems. They have attracted much research interest because of the various ecological functions of the mangrove ecosystems, including runoff and flood prevention, storage and recycling of nutrients and wastes, cultivation and energy conversion. The forests are major blue carbon systems, storing considerable amounts of carbon in marine sediments, thus becoming important regulators of climate change. Marine microorganisms are key parts of these mangrove ecosystems. However, much remains to be discovered about how mangrove microbiomes contribute to high ecosystem productivity and efficient cycling of elements.
Overview
There are about 80 different species of mangrove trees. All of these trees grow in areas with low-oxygen soil, where slow-moving waters allow fine sediments to accumulate. Mangrove forests grow only at tropical and subtropical latitudes near the equator because they cannot withstand freezing temperatures. Many mangrove forests can be recognised by their dense tangle of prop roots that make the trees appear to be standing on stilts above the water. This tangle of roots allows the trees to handle the daily rise and fall of tides, which means that most mangroves get flooded at least twice per day. The roots slow the movement of tidal waters, causing sediments to settle out of the water and build up the muddy bottom. Mangrove forests stabilise the coastline, reducing erosion from storm surges, currents, waves, and tides. The intricate root system of mangroves makes these forests attractive to fishes and other organisms seeking food and shelter from predators.
The main contribution of mangroves to the larger ecosystem comes from litter fall from the trees, which is then decomposed by primary consumers. Bacteria and protozoans colonise the plant litter and break it down chemically into organic compounds, minerals, carbon dioxide, and nitrogenous wastes. The intertidal existence to which these trees are adapted represents the major limitation to the number of species able to thrive in their habitat. High tide brings in salt water, and when the tide recedes, solar evaporation of the seawater in the soil leads to further increases in salinity. The return of tide can flush out these soils, bringing them back to salinity levels comparable to that of seawater. At low tide, organisms are exposed to increases in temperature and reduced moisture before being then cooled and flooded by the tide. Thus, for a plant to survive in this environment, it must tolerate broad ranges of salinity, temperature, and moisture, as well as several other key environmental factors—thus only a select few species make up the mangrove tree community.
A mangrove swamp typically features only a small number of tree species. It is not uncommon for a mangrove forest in the Caribbean to feature only three or four tree species. For comparison, a tropical rainforest biome may contain thousands of tree species, but this is not to say mangrove forests lack diversity. Though the trees are few in species, the ecosystem that these trees create provides a habitat for a great variety of other species, including as many as 174 species of marine megafauna.
Mangrove plants require a number of physiological adaptations to overcome the problems of low environmental oxygen levels, high salinity, and frequent tidal flooding. Each species has its own solutions to these problems; this may be the primary reason why, on some shorelines, mangrove tree species show distinct zonation. Small environmental variations within a mangal may lead to greatly differing methods for coping with the environment. Therefore, the mix of species is partly determined by the tolerances of individual species to physical conditions, such as tidal flooding and salinity, but may also be influenced by other factors, such as crabs preying on plant seedlings.
Once established, mangrove roots provide an oyster habitat and slow water flow, thereby enhancing sediment deposition in areas where it is already occurring. The fine, anoxic sediments under mangroves act as sinks for a variety of heavy (trace) metals which colloidal particles in the sediments have concentrated from the water. Mangrove removal disturbs these underlying sediments, often creating problems of trace metal contamination of seawater and organisms of the area.
Mangrove swamps protect coastal areas from erosion, storm surge (especially during tropical cyclones), and tsunamis. They limit high-energy wave erosion mainly during events such as storm surges and tsunamis. The mangroves' massive root systems are efficient at dissipating wave energy. Likewise, they slow down tidal water enough so that its sediment is deposited as the tide comes in, leaving all except fine particles when the tide ebbs. In this way, mangroves build their environments. Because of the uniqueness of mangrove ecosystems and the protection against erosion they provide, they are often the object of conservation programs, including national biodiversity action plans.
Distribution
Worldwide there are about 80 described species of mangroves that live along marine coasts. About 60 of these species are true mangroves which live only in the intertidal zone between high and low tides. "Mangroves once covered three-quarters of the world's tropical coastlines, with Southeast Asia hosting the greatest diversity. Only 12 species live in the Americas. Mangroves range in size from small bushes to the 60-meter giants found in Ecuador. Within a given mangrove forest, different species occupy distinct niches. Those that can handle tidal soakings grow in the open sea, in sheltered bays, and on fringe islands. Trees adapted to drier, less salty soil can be found farther from the shoreline. Some mangroves flourish along riverbanks far inland, as long as the freshwater current is met by ocean tides."
Mangroves can be found in 118 countries and territories in the tropical and subtropical regions of the world.
The largest percentage of mangroves is found between the 5° N and 5° S latitudes. Approximately 75% of world's mangroves are found in just 15 countries. Estimates of mangrove area based on remote sensing and global data tend to be lower than estimates based on literature and surveys for comparable periods.
In 2018, the Global Mangrove Watch Initiative released a global baseline based on remote sensing and global data for 2010. They estimated the total mangrove forest area of the world as of 2010 at , spanning 118 countries and territories. Following the conventions for identifying geographic regions from the Ramsar Convention on Wetlands, researchers reported that Asia has the largest share (38.7%) of the world's mangroves, followed by Latin America and the Caribbean (20.3%), Africa (20.0%), Oceania (11.9%), and Northern America (8.4%).
Sundarbans
The largest mangrove forest in the world is in the Sundarbans. The Sundarban forest lies in the vast delta on the Bay of Bengal formed by the super confluence of the Brahmaputra and Meghna rivers with distributaries of the Ganges. The seasonally flooded Sundarbans freshwater swamp forests lie inland from the mangrove forests on the coastal fringe. The forest covers of which about are in Bangladesh.
The Sundarbans is intersected by a complex network of tidal waterways, mudflats and small islands of salt-tolerant mangrove forests. The interconnected network of waterways makes almost every portion of the forest accessible by boat. The area is known as an important habitat for the endangered Bengal tiger, as well as numerous fauna including species of birds, spotted deer, crocodiles and snakes. The fertile soils of the delta have been subject to intensive human use for centuries, and the ecoregion has been mostly converted to intensive agriculture, with few enclaves of forest remaining. Additionally, the Sundarbans serves a crucial function as a protective barrier for millions of inhabitants against floods that result from cyclones.
Four protected areas in the Sundarbans are listed as UNESCO World Heritage Sites. Despite these protections, the Indian Sundarbans were assessed as endangered in 2020 under the IUCN Red List of Ecosystems framework. There is a consistent pattern of depleted biodiversity or loss of species and the ecological quality of the forest is declining.
Ecosystem
The unique ecosystem found in the intricate mesh of mangrove roots offers a quiet marine habitat for young organisms. In areas where roots are permanently submerged, the organisms they host include algae, barnacles, oysters, sponges, and bryozoa, which all require a hard surface for anchoring while they filter-feed. Shrimp and mud lobsters use the muddy bottoms as their home. Mangrove crabs eat the mangrove leaves, adding nutrients to the mangal mud for other bottom feeders. In at least some cases, the export of carbon fixed in mangroves is important in coastal food webs. Mangrove plantations host several commercially important species of fish and crustaceans.
In Puerto Rico, the red, white, and black mangroves occupy different ecological niches and have slightly different chemical compositions, so the carbon content varies between the species, as well between the different tissues of the plant (e.g., leaf matter versus roots). There is a clear succession of these three trees from the lower elevations, which are dominated by red mangroves, to farther inland with a higher concentration of white mangroves.
Mangrove forests are an important part of the cycling and storage of carbon in tropical coastal ecosystems. Knowing this, scientists seek to reconstruct the environment and investigate changes to the coastal ecosystem over thousands of years using sediment cores. However, an additional complication is the imported marine organic matter that also gets deposited in the sediment through the tidal flushing of mangrove forests.
Mangrove forests can decay into peat deposits because of fungal and bacterial processes as well as by the action of termites. It becomes peat in good geochemical, sedimentary, and tectonic conditions. The nature of these deposits depends on the environment and the types of mangroves involved. Termites process fallen leaf litter, root systems and wood from mangroves into peat to build their nests. Termites stabilise the chemistry of this peat and represent approximately 2% of above ground carbon storage in mangroves. As the nests are buried over time this carbon is stored in the sediment, and the carbon cycle continues.
Mangroves are an important source of blue carbon. Globally, mangroves stored of carbon in 2012. Two percent of global mangrove carbon was lost between 2000 and 2012, equivalent to a maximum potential of of CO2 emissions. Globally, mangroves have been shown to provide measurable economic protections to coastal communities affected by tropical storms.
Biodiversity
Birds
Heterogeneity in landscape ecology is a measure of how different parts of a landscape are from one another. It can manifest in an ecosystem from the abiotic or biotic characteristics of the environment. For example, coastal mangrove forests are located at the land-sea interface, so their functioning is influenced by abiotic factors such as tides, as well as biotic factors such as the extent and configuration of adjacent vegetation. For forest birds, tidal inundation means that the availability of many mangrove resources fluctuates daily, suggesting foraging flexibility is likely to be important. Mangroves also offer estuarine prey items, such as mudskippers and crabs, that are not found in terrestrial forest types. Further, mangroves are often situated in a complex mosaic of adjacent vegetation types such as grasslands, saltmarshes, and woodlands, and this can mean that flexibility in foraging strategy and choice of foraging habitat may be advantageous for highly mobile forest birds. Relative to other forest types, mangroves support few bird species that are obligate habitat (mangrove) specialists and instead host many species with generalised foraging niches.
Bird sanctuaries
Mangrove forests are home and sanctuaries for many of aquatic bird species, including:
Pulicat Lake Bird Sanctuary
Mangalavanam Bird Sanctuary
Salim Ali Bird Sanctuary
Pichavaram
Coringa Wildlife Sanctuary
Sundarbans National Park
Fish
The intricate root system of mangrove forests makes them attractive to adult fish seeking food and juvenile fish seeking shelter.
Mangrove crabs and holobionts
Mangrove forests are among the more productive and diverse ecosystems on the planet, despite limited nitrogen availability. Under such conditions, animal-microbe associations (holobionts) are often key to ecosystem functioning. An example is the role of fiddler crabs and their carapace-associated microbial biofilm as hotspots of microbial nitrogen transformations and sources of nitrogen within the mangrove ecosystem.
Among coastal ecosystems, mangrove forests are of great importance as they account for three quarters of the tropical coastline and provide different ecosystem services. Mangrove ecosystems generally act as a net sink of carbon, although they release organic matter to the sea in the form of dissolved refractory macromolecules, leaves, branches and other debris. In pristine environments, mangroves are among the most productive ecosystems on the planet, despite growing in tropical waters that are often nutrient depleted. The refractory nature of the organic matter produced and retained in mangroves can slow the recycling of nutrients, particularly of nitrogen. Nitrogen limitation in such systems may be overcome by microbial nitrogen fixation when combined with high rates of bioturbation by macrofauna, such as crabs and lobsters.
Bioturbation by macrofauna affect nitrogen availability and multiple nitrogen related microbial processes through sediment reworking, burrow construction and bioirrigation, feeding and excretion. Macrofauna mix old and fresh organic matter, extend oxic–anoxic sediment interfaces, increase the availability of energy-yielding electron acceptors and increase nitrogen turnover via direct excretion. Thus, macrofauna may alleviate nitrogen limitation by priming the remineralisation of refractory nitrogen (that is, the nitrogen that can't be biologically decomposed), reducing plant-microbe competition. Such activity ultimately promotes nitrogen recycling, plant assimilation and high nitrogen retention, as well as favours its loss by stimulating coupled nitrification and denitrification.
Mangrove sediments are highly bioturbated by decapods such as crabs. Crab populations continuously rework sediment by constructing burrows, creating new niches, transporting or selectively grazing on sediment microbial communities. In addition, crabs can affect organic matter turnover by assimilating leaves and producing finely fragmented faeces, or by carrying them into their burrows. Therefore, crabs are considered important ecosystem engineers shaping biogeochemical processes in intertidal muddy banks of mangroves. In contrast to burrowing polychaetes or amphipods, the abundant Ocipodid crabs, mainly represented by fiddler crabs, do not permanently ventilate their burrows. These crabs may temporarily leave their burrows for surface activities, or otherwise plug their burrow entrance during tidal inundation in order to trap air. A recent study showed that these crabs can be associated with a diverse microbial community, either on their carapace or in their gut.
The exoskeleton of living animals, such as shells or carapaces, offers a habitat for microbial biofilms which are actively involved in different N-cycling pathways such as nitrification, denitrification and dissimilatory nitrate reduction to ammonium (DNRA). Colonizing the carapace of crabs may be advantageous for specific bacteria, because of host activities such as respiration, excretion, feeding and horizontal and vertical migrations. However, the ecological interactions between fiddler crabs and bacteria, their regulation and significance as well as their implications at scales spanning from the single individual to the ecosystem are not well understood.
Biogeochemistry
Carbon cycle
Mangrove forests are amongst the world's most productive marine ecosystems, with net primary productivity (NPP) in the order of 208 Tg C yr−1. Mangrove forests achieve a steady state once the forest reaches maximum biomass at around 20–30 years through a constant process of mortality and renewal so, assuming the living biomass is not becoming more carbon dense, then carbon has to be lost at a rate equal to the amount of carbon fixed as NPP. Hence this productivity is either retained within the mangrove forest, as a standing stock of live material such as wood, buried in sediments, or exported to neighbouring habitats as litter, particulate and dissolved organic carbon (POC and DOC) and dissolved inorganic carbon (DIC), or lost to the atmosphere.
The outwelling hypothesis argues that export of locally derived POC and DOC is an important ecosystem function of mangroves, which drives detrital based food webs in adjacent coastal habitats. Export of mangrove carbon has been estimated to make a significant trophic contribution to adjacent ecosystems. The theory of outwelling is supported by mass balance evaluations that show the amount of carbon fixed by mangroves normally greatly exceeds the amount stored within the forest, although the scale of outwelling varies considerably between forests, due to differences in coastal geomorphology, tidal regimes, freshwater flow and productivity.
In the 1990s, global estimates could account for 48% of the total global mangrove primary production of 218 ± 72 million tons C yr−1 (see diagram on the right). By incorporating information on carbon burial, efflux and carbon outwelled as leaf litter, POC and DOC, the remaining 52% was thought outwelled as DIC, though there was insufficient data to confirm this. More recent assessments of DIC export at two sites in Australia supported the estimates of Bouillon et al. in 2008, although in 2014 Alongi suggested that only 40% of NPP was exported as DIC.
Nitrogen assimilation
Mangrove forests and coastal marshes are typically considered N-limited ecosystems because of their high primary production. Therefore, mangrove plants are highly efficient at utilising soil nitrogen, making them an important sink for excess nitrogen from upstream. However, different mangrove species may still utilise nitrogen at different efficiencies, even though they share similar nitrogen pathways (see diagram on right). Reported nitrogen assimilation rates in mangrove plants ranged from 2 to 8 μmol g−1 h−1 under ambient nitrogen conditions, and 19 to 251 μmol g−1 h−1 when the nitrogen supply was unlimited.
In addition to species variation, different environmental conditions can also affect the nitrogen assimilation rates in mangrove plants. Because Cl− ions can reduce protein synthesis and nitrogen assimilation, soil pore water salinity appears to be a negative factor that significantly alters the nitrogen uptake rates of mangrove plants.
Exploitation and conservation
Adequate data is only available for about half of the global area of mangroves. However, of those areas for which data has been collected, it appears that 35% of the mangroves have been destroyed. Since the 1980s, around 2% of mangrove area is estimated to be lost each year. Assessments of global variation in mangrove loss indicates that national regulatory quality mediates how different drivers and pressures influence loss rates.
Shrimp farming causes approximately a quarter of the destruction of mangrove forests. Likewise, the 2010 update of the World Mangrove Atlas indicated that approximately one fifth of the world's mangrove ecosystems have been lost since 1980, although this rapid loss rate appears to have decreased since 2000 with global losses estimated at between 0.16% and 0.39% annually between 2000 and 2012. Despite global loss rates decreasing since 2000, Southeast Asia remains an area of concern with loss rates between 3.6% and 8.1% between 2000 and 2012. By far the most damaging form of shrimp farming is when a closed ponds system (non-integrated multi-trophic aquaculture) is used, as these require destruction of a large part of the mangrove, and use antibiotics and disinfectants to suppress diseases that occur in this system, and which may also leak into the surrounding environment. Far less damage occurs when integrated mangrove-shrimp aquaculture is used, as this is connected to the sea and subjected to the tides, and less diseases occur, and as far less mangrove is destroyed for it.
Grassroots efforts to protect mangroves from development and from citizens cutting down the mangroves for charcoal production, cooking, heating and as a building material are becoming more popular. Solar cookers are distributed by many non-government organizations as a low-cost alternative to wood and charcoal stoves. These may help in reducing the demand for charcoal.
In Thailand, community management has been effective in restoring damaged mangroves. Also, production of mangrove honey is practiced, as a way to generate sustainable income for nearby people, keeping them from destroying the mangrove and generate a short-term revenue.
In Madagascar, honey is also produced in mangroves as a source of (non-destructive) income generation. In addition, silk pods from endemic silkworm species are also collected in the Madagascar mangroves for wild silk production.
In the Bahamas, for example, active efforts to save mangroves are occurring on the islands of Bimini and Great Guana Cay.
In Trinidad and Tobago as well, efforts are underway to protect a mangrove threatened by the construction of a steel mill and a port.
Within northern Ecuador, mangrove regrowth is reported in almost all estuaries and stems primarily from local actors responding to earlier periods of deforestation in the Esmeraldas region.
Mangroves have been reported to be able to help buffer against tsunami, cyclones, and other storms, and as such may be considered a flagship system for ecosystem-based adaptation to the impacts of climate change. One village in Tamil Nadu was protected from tsunami destruction—the villagers in Naluvedapathy planted 80,244 saplings to get into the Guinness Book of World Records. This created a kilometre-wide belt of trees of various varieties. When the 2004 tsunami struck, much of the land around the village was flooded, but the village suffered minimal damage.
Ocean deoxygenation
Compared to seagrass meadows and coral reefs, hypoxia is more common on a regular basis in mangrove ecosystems, through ocean deoxygenation is compounding the negative effects by anthropogenic nutrient inputs and land use modification.
Like seagrass, mangrove trees transport oxygen to roots of rhizomes, reduce sulfide concentrations, and alter microbial communities. Dissolved oxygen is more readily consumed in the interior of the mangrove forest. Anthropogenic inputs may push the limits of survival in many mangrove microhabitats. For example, shrimp ponds constructed in mangrove forests are considered the greatest anthropogenic threat to mangrove ecosystems. These shrimp ponds reduce estuary circulation and water quality which leads to the promotion of diel-cycling hypoxia. When the quality of the water degrades, the shrimp ponds are quickly abandoned leaving massive amounts of wastewater. This is a major source of water pollution that promotes ocean deoxygenation in the adjacent habitats.
Due to these frequent hypoxic conditions, the water does not provide habitats to fish. When exposed to extreme hypoxia, ecosystem function can completely collapse. Extreme deoxygenation will affect the local fish populations, which are an essential food source. The environmental costs of shrimp farms in the mangrove forests grossly outweigh the economic benefits of them. Cessation of shrimp production and restoration of these areas reduce eutrophication and anthropogenic hypoxia.
Reforestation
In some areas, mangrove reforestation and mangrove restoration is also underway. Red mangroves are the most common choice for cultivation, used particularly in marine aquariums in a sump to reduce nitrates and other nutrients in the water. Mangroves also appear in home aquariums, and as ornamental plants, such as in Japan.
The Manzanar Mangrove Initiative is an ongoing experiment in Arkiko, Eritrea, part of the Manzanar Project founded by Gordon H. Sato, establishing new mangrove plantations on the coastal mudflats. Initial plantings failed, but observation of the areas where mangroves did survive by themselves led to the conclusion that nutrients in water flow from inland were important to the health of the mangroves. Trials with the Eritrean Ministry of Fisheries followed, and a planting system was designed to provide the nitrogen, phosphorus, and iron missing from seawater.
The propagules are planted inside a reused galvanized steel can with the bottom knocked out; a small piece of iron and a pierced plastic bag with fertilizer containing nitrogen and phosphorus are buried with the propagule. , after six years of planting, 700,000 mangroves are growing; providing stock feed for sheep and habitat for oysters, crabs, other bivalves, and fish.
Another method of restoring mangroves is by using quadcopters (which are able to carry and deposit seed pods). According to Irina Fedorenko, an amount of work equivalent to weeks of planting using traditional methods can be done by a drone in days, and at a fraction of the cost.
Seventy percent of mangrove forests have been lost in Java, Indonesia. Mangroves formerly protected the island's coastal land from flooding and erosion. Wetlands International, an NGC based in the Netherlands, in collaboration with nine villages in Demak where lands and homes had been flooded, began reviving mangrove forests in Java. Wetlands International introduced the idea of developing tropical versions of techniques traditionally used by the Dutch to catch sediment in North Sea coastal salt marshes. Originally, the villagers constructed a sea barrier by hammering two rows of vertical bamboo poles into the seabed and filling the gaps with brushwood held in place with netting. Later the bamboo was replaced by PVC pipes filled with concrete. As sediment gets deposited around the brushwood, it serves to catch floating mangrove seeds and provide them with a stable base to germinate, take root and regrow. This creates a green belt of protection around the islands. As the mangroves mature, more sediment is held in the catchment area; the process is repeated until a mangrove forest has been restored. Eventually the protective structures will not be needed. By late 2018, of brushwood barriers along the coastline had been completed.
A concern over reforestation is that although it supports increases in mangrove area it may actually result in a decrease in global mangrove functionality and poor restoration processes may result in longer term depletion of the mangrove resource.
National and international studies
In terms of local and national studies of mangrove loss, the case of Belize's mangroves is illustrative in its contrast to the global picture. A recent, satellite-based study—funded by the World Wildlife Fund and conducted by the Water Center for the Humid Tropics of Latin America and the Caribbean (CATHALAC)—indicates Belize's mangrove cover declined by a mere 2% over a 30-year period. The study was born out of the need to verify the popular conception that mangrove clearing in Belize was rampant.
Instead, the assessment showed, between 1980 and 2010, under of mangroves had been cleared, although clearing of mangroves near Belize's main coastal settlements (e.g. Belize City and San Pedro) was relatively high. The rate of loss of Belize's mangroves—at 0.07% per year between 1980 and 2010—was much lower than Belize's overall rate of forest clearing (0.6% per year in the same period). These findings can also be interpreted to indicate Belize's mangrove regulations (under the nation's) have largely been effective. Nevertheless, the need to protect Belize's mangroves is imperative, as a 2009 study by the World Resources Institute (WRI) indicates the ecosystems contribute to million per year to Belize's national economy.
From 1990, in Tanzania, Adelaida K. Semesi led aresearch programme which resulted in Tanzania being one of the first countries to have an environmental management plan for mangroves. Nicknamed "mama mikoko" ("mama mangroves" in Swahili), Semesi also was a Council Member for the International Society for Mangrove Ecosystems.
In May 2019, ORNL DAAC News announced that NASA's Carbon Monitoring System (CMS), using new satellite-based maps of global mangrove forests across 116 countries, had created a new dataset to characterize the "distribution, biomass, and canopy height of mangrove-forested wetlands". Mangrove forests move carbon dioxide "from the atmosphere into long-term storage" in greater quantities than other forests, making them "among the planet's best carbon scrubbers" according to a NASA-led study.
| Physical sciences | Oceanography | Earth science |
2611797 | https://en.wikipedia.org/wiki/Moschops | Moschops | Moschops (Ancient Greek for "calf face") is an extinct genus of therapsids that lived in the Guadalupian epoch, around 265–260 million years ago. They were heavily built plant eaters, and they may have lived partly in water, as hippopotamuses do. They had short, thick heads and might have competed by head-butting each other. Their elbow joints allowed them to walk with a more mammal-like gait rather than crawling. Their remains were found in the Karoo region of South Africa, belonging to the Tapinocephalus Assemblage Zone. Therapsids, such as Moschops, are synapsids, the dominant land animals in the Permian period, which ended 252 million years ago.
Description
Moschops were heavy set dinocephalian synapsids, measuring in length, and weighing on average and in maximum body mass. They had small heads with broad orbits and heavily built short necks. Like other members of Tapinocephalidae, the skull had a tiny opening for the pineal organ. The occiput was broad and deep, but the skull was more narrow in the dorsal border. Furthermore, the pterygoid arches and the angular region of the jaw with heavily built jaw muscles. Due to that and the possession of long-crowned, stout teeth, it is believed that Moschops was a herbivore feeding on nutrient-poor and tough vegetation, like cycad stems. Due to the presumably nutrient-poor food, it is likely they had to feed for long periods of time. The anatomy of the taxa allowed them to open the elbow joints more widely, enabling them to move in a more mammal-like posture than some other animals at the time. This helped to carry their massive bodies more easily while feeding, as well as allowing them short bursts of speed. It has also been proposed that Moschops were possibly sub-aquatic. Moschops had rather thick skulls, prompting speculation that individuals could have competed with one another by head-butting. A 2017 published study would later confirm this by synchrotron scanning a Moschops capensis skull, which revealed numerous anatomical adaptations to the central nervous system for combative behaviour. They were likely preyed upon by titanosuchids and larger therocephalian species.
Earliest finds
Moschops material was first discovered in the Ecca Group (part of the Karoo Supergroup) of South Africa by Robert Broom. As the geological horizon was dubious, it was referred to have originated from the Ecca Group on the basis of Pareiasaurus remains in near proximity. The discovered material includes a holotype (AMNH 5550) and seven topotypes (AMNH 5551-5557). The degree of pachyostosis varies within the skulls of the specimens, and Broom believed this to have been linked to variations in gender and age. In 1910, the material was sent to the American Museum of Natural History in New York City and described in 1911.
Classification
Moschops is characterized by a strongly pachyostosed skull with a broad intertemporal region and greatly reduced temporal fossae. Two species are known from the fossil record, M. capensis and M. koupensis. Two other species were assigned (M. whaitsi and M. oweni), but their validity is considered possibly dubious. Genera regarded as synonyms are Moschoides, Agnosaurus, Moschognathus and Pnigalion. Delphinognathus conocephalus could represent juvenile Moschops, thus possibly synonymous. Delphinognathus is only known from a single, moderately pachyostosed skull. It has a conical boss on the parietal surrounding the pineal foramen.
| Biology and health sciences | Proto-mammals | Animals |
2612628 | https://en.wikipedia.org/wiki/Retractor%20%28medicine%29 | Retractor (medicine) | A retractor is a surgical instrument used to separate the edges of a surgical incision/wound or to hold away certain organs and tissues (i.e. to provide tissue retraction) so that body parts underneath may be accessed during surgical operations.
The broad term retractor typically describes a simple steel tool possessing a curved, hooked, or angled blade, which is manually manipulated to help maintain a desired position of tissue during surgery. More sophisticated retractors may be clamped in place (usually to a tableside frame) or suspended at the end of a robotic arm. Retractors can also be "self-retaining" and no longer need to be held once inserted, having two or more opposing blades or hooks which are separated via spring, ratchet, worm gear or other method. The term retractor is also used to describe distinct, hand-cranked devices such as rib spreaders (also known as thoracic retractors, or distractors) with which surgeons may use to forcefully drive tissues apart to obtain exposure. Different surgery specialties can have specific kinds of retractors - e.g., for certain kinds of spinal surgery, such as Minimally Invasive Transforaminal Lumbar Interbody Fusions, some retractors are fitted both with suction and with fiberoptic lights to keep deep surgical wounds both dry and illuminated.
Surgical assistants, whether they be another surgeon, surgical residents or professionally trained procedure assistants (specifically Certified Surgical Assistants, Registered Nurse First Assistants, Physicians Assistants, or Surgical Technologists), may assist the operating surgeon in the process of retraction.
History
Surgical retractors probably originate with very basic tools used in the Stone Age. Branches or antlers of various shapes were used to dig and extract food from the ground. As the use of tools evolved, a variety of instruments came about to substitute for the use of hooked or grasping fingers in the butchering of meat or dissection of bodies. The use of metals in tool making was of great importance. A variety of Roman metal instruments of the hook and retractor family have been found by archeologists. These instruments would generally be called hooks if the end was as narrow as the handle of the instrument. If the end was broad, it would be called a retractor. Also arising from this group of tools were other related tools for displacing (elevators and spatulas) and scooping (spoons and curettes).
In 4th century CE, Indian physician Susruta used surgical tools such as retractors. In a description of the procedure of tonsillectomy from the 7th century CE, Paul of Aegina documents the use of a tongue spatula to keep the tongue out of the way while a form of tonsil hook is used to bring the tonsil forward for excision. In 1000 CE Abu al-Qasim al-Zahrawi, also known as Albucasis or Abulcasis, described a variety of surgical instruments including retractors in his famous text Al-Tasrif. Vesalius described a variety of hooks and retractors in the 16th century. Jan Mikulicz-Radecki's invention of a hinged rib spreading retractor in 1904 prompted a flurry of development of retractors in the early 20th century, culminating in 1936 in our modern device based on the design of Enrique Finochietto.
Current
The following is an incomplete list of surgical retractors in use:
In 2021, a self-retaining ring retractor in the UK was given HRH The Queen's Award for Enterprise in the Innovation category as the first and only retractor ever to be awarded.
Hand-held retractors
Hohmann Retractor
Lahey Retractor
Senn Retractor
Blair (Rollet) Retractor
Rigid Rake
Flexible Rake
"Cat's Paws" - sharp, wide rakes
Ragnell Retractor
Linde-Ragnell Retractor
Davis Retractor
Volkman Retractor
Farabeuf Retractor
Mathieu Retractor
Jackson Tracheal Hook
Meyerding Finger Retractor
Little Retractor
Love Nerve Retractor
Green Retractor
Goelet Retractor
Cushing Vein Retractor
Langenbeck Retractor
Richardson Retractor
Richardson-Eastmann Retractor
Deaver Retractor
Doyen Retractor
Parker Retractor
Parker-Mott Retractor
Roux Retractor
Mayo-Collins Retractor
"Army-Navy" Retractor
Ribbon Retractor - malleable, able to be bent as the surgeon desires
Self-retaining retractors
Rultract Skyhook Retractor System
Alms Retractor
Lone Star Retractor
Galaxy II retractor
Gelpi Retractor
Gutow Retractor
Weitlaner Retractor
Beckman-Weitlaner Retractor
Beckman-Eaton Retractor
Beckman Retractor
Balfour Retractor - typically used in lower abdomen and pelvic surgery
Rib spreader (a.k.a. Finochietto retractor)
Travers Retractor
West Retractor
Norfolk & Norwich Retractor
Gallery
| Technology | Surgical instruments | null |
6259941 | https://en.wikipedia.org/wiki/Biomaterial | Biomaterial | A biomaterial is a substance that has been engineered to interact with biological systems for a medical purpose – either a therapeutic (treat, augment, repair, or replace a tissue function of the body) or a diagnostic one. The corresponding field of study, called biomaterials science or biomaterials engineering, is about fifty years old. It has experienced steady growth over its history, with many companies investing large amounts of money into the development of new products. Biomaterials science encompasses elements of medicine, biology, chemistry, tissue engineering and materials science.
A biomaterial is different from a biological material, such as bone, that is produced by a biological system. However, "biomaterial" and "biological material" are often used interchangeably. Further, the word "bioterial" has been proposed as a potential alternate word for biologically-produced materials such as bone, or fungal biocomposites. Additionally, care should be exercised in defining a biomaterial as biocompatible, since it is application-specific. A biomaterial that is biocompatible or suitable for one application may not be biocompatible in another.
Introduction
Biomaterials can be derived either from nature or synthesized in the laboratory using a variety of chemical approaches utilizing metallic components, polymers, ceramics or composite materials. They are often used and/or adapted for a medical application, and thus comprise the whole or part of a living structure or biomedical device which performs, augments, or replaces a natural function. Such functions may be relatively passive, like being used for a heart valve, or maybe bioactive with a more interactive functionality such as hydroxy-apatite coated hip implants. Biomaterials are also commonly used in dental applications, surgery, and drug delivery. For example, a construct with impregnated pharmaceutical products can be placed into the body, which permits the prolonged release of a drug over an extended period of time. A biomaterial may also be an autograft, allograft or xenograft used as a transplant material.
Bioactivity
The ability of an engineered biomaterial to induce a physiological response that is supportive of the biomaterial's function and performance is known as bioactivity. Most commonly, in bioactive glasses and bioactive ceramics this term refers to the ability of implanted materials to bond well with surrounding tissue in either osteo conductive or osseo productive roles. Bone implant materials are often designed to promote bone growth while dissolving into surrounding body fluid. Thus for many biomaterials good biocompatibility along with good strength and dissolution rates are desirable. Commonly, bioactivity of biomaterials is gauged by the surface biomineralization in which a native layer of hydroxyapatite is formed at the surface. These days, the development of clinically useful biomaterials is greatly enhanced by the advent of computational routines that can predict the molecular effects of biomaterials in a therapeutic setting based on limited in vitro experimentation.
Self-assembly
Self-assembly is the most common term in use in the modern scientific community to describe the spontaneous aggregation of particles (atoms, molecules, colloids, micelles, etc.) without the influence of any external forces. Large groups of such particles are known to assemble themselves into thermodynamically stable, structurally well-defined arrays, quite reminiscent of one of the seven crystal systems found in metallurgy and mineralogy (e.g., face-centered cubic, body-centered cubic, etc.). The fundamental difference in equilibrium structure is in the spatial scale of the unit cell (lattice parameter) in each particular case.
Molecular self assembly is found widely in biological systems and provides the basis of a wide variety of complex biological structures. This includes an emerging class of mechanically superior biomaterials based on microstructural features and designs found in nature. Thus, self-assembly is also emerging as a new strategy in chemical synthesis and nanotechnology. Molecular crystals, liquid crystals, colloids, micelles, emulsions, phase-separated polymers, thin films and self-assembled monolayers all represent examples of the types of highly ordered structures, which are obtained using these techniques. The distinguishing feature of these methods is self-organization.
Structural hierarchy
Nearly all materials could be seen as hierarchically structured, since the changes in spatial scale bring about different mechanisms of deformation and damage. However, in biological materials, this hierarchical organization is inherent to the microstructure. One of the first examples of this, in the history of structural biology, is the early X-ray scattering work on the hierarchical structure of hair and wool by Astbury and Woods. In bone, for example, collagen is the building block of the organic matrix, a triple helix with diameter of 1.5 nm. These tropocollagen molecules are intercalated with the mineral phase (hydroxyapatite, calcium phosphate) forming fibrils that curl into helicoids of alternating directions. These "osteons" are the basic building blocks of bones, with the volume fraction distribution between organic and mineral phase being about 60/40.
In another level of complexity, the hydroxyapatite crystals are mineral platelets that have a diameter of approximately 70 to 100 nm and thickness of 1 nm. They originally nucleate at the gaps between collagen fibrils.
Similarly, the hierarchy of abalone shell begins at the nanolevel, with an organic layer having a thickness of 20 to 30 nm. This layer proceeds with single crystals of aragonite (a polymorph of CaCO3) consisting of "bricks" with dimensions of 0.5 and finishing with layers approximately 0.3 mm (mesostructure).
Crabs are arthropods, whose carapace is made of a mineralized hard component (exhibits brittle fracture) and a softer organic component composed primarily of chitin. The brittle component is arranged in a helical pattern. Each of these mineral "rods" (1 μm diameter) contains chitin–protein fibrils with approximately 60 nm diameter. These fibrils are made of 3 nm diameter canals that link the interior and exterior of the shell.
Applications
Biomaterials are used in:
Joint replacements
Bone plates
Intraocular lenses (IOLs) for eye surgery
Bone cement
Artificial ligaments and tendons
Dental implants for tooth fixation
Blood vessel prostheses
Heart valves
Skin repair devices (artificial tissue)
Cochlear replacements
Contact lenses
Breast implants
Drug delivery mechanisms
Sustainable materials
Vascular grafts
Stents
Nerve conduits
Surgical sutures, clips, and staples for wound closure
Pins and screws for fracture stabilisation
Surgical mesh
Biomaterials must be compatible with the body, and there are often issues of biocompatibility, which must be resolved before a product can be placed on the market and used in a clinical setting. Because of this, biomaterials are usually subjected to the same requirements as those undergone by new drug therapies. All manufacturing companies are also required to ensure traceability of all of their products, so that if a defective product is discovered, others in the same batch may be traced.
Bone grafts
Calcium sulfate (its α- and β-hemihydrates) is a well known biocompatible material that is widely used as a bone graft substitute in dentistry or as its binder.
Heart valves
In the United States, 49% of the 250,000 valve replacement procedures performed annually involve a mechanical valve implant. The most widely used valve is a bileaflet disc heart valve or St. Jude valve. The mechanics involve two semicircular discs moving back and forth, with both allowing the flow of blood as well as the ability to form a seal against backflow. The valve is coated with pyrolytic carbon and secured to the surrounding tissue with a mesh of woven fabric called Dacron (du Pont's trade name for polyethylene terephthalate). The mesh allows for the body's tissue to grow, while incorporating the valve.
Skin repair
Most of the time, artificial tissue is grown from the patient's own cells. However, when the damage is so extreme that it is impossible to use the patient's own cells, artificial tissue cells are grown. The difficulty is in finding a scaffold that the cells can grow and organize on. The characteristics of the scaffold must be that it is biocompatible, cells can adhere to the scaffold, mechanically strong and biodegradable. One successful scaffold is a copolymer of lactic acid and glycolic acid.
Properties
As discussed previously, biomaterials are used in medical devices to treat, assist, or replace a function within the human body. The application of a specific biomaterial must combine the necessary composition, material properties, structure, and desired in vivo reaction in order to perform the desired function. Categorizations of different desired properties are defined in order to maximize functional results.
Host response
Host response is defined as the "response of the host organism (local and systemic) to the implanted material or device". Most materials will have a reaction when in contact with the human body. The success of a biomaterial relies on the host tissue's reaction with the foreign material. Specific reactions between the host tissue and the biomaterial can be generated through the biocompatibility of the material.
Biomaterial and tissue interactions
The in vivo functionality and longevity of any implantable medical device is affected by the body's response to the foreign material. The body undergoes a cascade of processes defined under the foreign body response (FBR) in order to protect the host from the foreign material. The interactions between the device upon the host tissue/blood as well as the host tissue/blood upon the device must be understood in order to prevent complications and device failure.
Tissue injury caused by device implantation causes inflammatory and healing responses during FBR. The inflammatory response occurs within two time periods: the acute phase, and the chronic phase. The acute phase occurs during the initial hours to days of implantation, and is identified by fluid and protein exudation along with a neutrophilic reaction. During the acute phase, the body attempts to clean and heal the wound by delivering excess blood, proteins, and monocytes are called to the site. Continued inflammation leads to the chronic phase, which can be categorized by the presence of monocytes, macrophages, and lymphocytes. In addition, blood vessels and connective tissue form in order to heal the wounded area.
Compatibility
Biocompatibility is related to the behavior of biomaterials in various environments under various chemical and physical conditions. The term may refer to specific properties of a material without specifying where or how the material is to be used. For example, a material may elicit little or no immune response in a given organism, and may or may not able to integrate with a particular cell type or tissue. Immuno-informed biomaterials that direct the immune response rather than attempting to circumvent the process is one approach that shows promise. The ambiguity of the term reflects the ongoing development of insights into "how biomaterials interact with the human body" and eventually "how those interactions determine the clinical success of a medical device (such as pacemaker or hip replacement)". Modern medical devices and prostheses are often made of more than one material, so it might not always be sufficient to talk about the biocompatibility of a specific material. Surgical implantation of a biomaterial into the body triggers an organism-inflammatory reaction with the associated healing of the damaged tissue. Depending upon the composition of the implanted material, the surface of the implant, the mechanism of fatigue, and chemical decomposition there are several other reactions possible. These can be local as well as systemic. These include immune response, foreign body reaction with the isolation of the implant with a vascular connective tissue, possible infection, and impact on the lifespan of the implant. Graft-versus-host disease is an auto- and alloimmune disorder, exhibiting a variable clinical course. It can manifest in either acute or chronic form, affecting multiple organs and tissues and causing serious complications in clinical practice, both during transplantation and implementation of biocompatible materials.
Toxicity
A biomaterial should perform its intended function within the living body without negatively affecting other bodily tissues and organs. In order to prevent unwanted organ and tissue interactions, biomaterials should be non-toxic. The toxicity of a biomaterial refers to the substances that are emitted from the biomaterial while in vivo. A biomaterial should not give off anything to its environment unless it is intended to do so. Nontoxicity means that biomaterial is: noncarcinogenic, nonpyrogenic, nonallergenic, blood compatible, and noninflammatory. However, a biomaterial can be designed to include toxicity for an intended purpose. For example, application of toxic biomaterial is studied during in vivo and in vitro cancer immunotherapy testing. Toxic biomaterials offer an opportunity to manipulate and control cancer cells. One recent study states: "Advanced nanobiomaterials, including liposomes, polymers, and silica, play a vital role in the codelivery of drugs and immunomodulators. These nanobiomaterial-based delivery systems could effectively promote antitumor immune responses and simultaneously reduce toxic adverse effects." This is a prime example of how the biocompatibility of a biomaterial can be altered to produce any desired function.
Biodegradable biomaterials
Biodegradable biomaterials refers to materials that are degradable through natural enzymatic reactions. The application of biodegradable synthetic polymers began in the later 1960s. Biodegradable materials have an advantage over other materials, as they have lower risk of harmful effects long term. In addition to ethical advancements using biodegradable materials, they also improve biocompatibility for materials used for implantation. Several properties including biocompatibility are important when considering different biodegradable biomaterials. Biodegradable biomaterials can be synthetic or natural depending on their source and type of extracellular matrix (ECM).
Biocompatible plastics
Some of the most commonly-used biocompatible materials (or biomaterials) are polymers due to their inherent flexibility and tunable mechanical properties. Medical devices made of plastics are often made of a select few including: cyclic olefin copolymer (COC), polycarbonate (PC), polyetherimide (PEI), medical grade polyvinylchloride (PVC), polyethersulfone (PES), polyethylene (PE), polyetheretherketone (PEEK) and even polypropylene (PP). To ensure biocompatibility, there are a series of regulated tests that material must pass to be certified for use. These include the United States Pharmacopoeia IV (USP Class IV) Biological Reactivity Test and the International Standards Organization 10993 (ISO 10993) Biological Evaluation of Medical Devices. The main objective of biocompatibility tests is to quantify the acute and chronic toxicity of material and determine any potential adverse effects during use conditions, thus the tests required for a given material are dependent on its end-use (i.e. blood, central nervous system, etc.).
Surface and bulk properties
Two properties that have a large effect on the functionality of a biomaterial is the surface and bulk properties.
Bulk properties refers to the physical and chemical properties that compose the biomaterial for its entire lifetime. They can be specifically generated to mimic the physiochemical properties of the tissue that the material is replacing. They are mechanical properties that are generated from a material's atomic and molecular construction.
Important bulk properties:
Chemical Composition
Microstructure
Elasticity
Tensile Strength
Density
Hardness
Electrical Conductivity
Thermal Conductivity
Surface properties refers to the chemical and topographical features on the surface of the biomaterial that will have direct interaction with the host blood/tissue. Surface engineering and modification allows clinicians to better control the interactions of a biomaterial with the host living system.
Important surface properties:
Wettability (surface energy)
Surface chemistry
Surface textures (smooth/rough)
Topographical factors including: size, shape, alignment, structure determine the roughness of a material.
Surface Tension
Surface Charge
Mechanical properties
In addition to a material being certified as biocompatible, biomaterials must be engineered specifically to their target application within a medical device. This is especially important in terms of mechanical properties which govern the way that a given biomaterial behaves. One of the most relevant material parameters is the Young's Modulus, E, which describes a material's elastic response to stresses. The Young's Moduli of the tissue and the device that is being coupled to it must closely match for optimal compatibility between device and body, whether the device is implanted or mounted externally. Matching the elastic modulus makes it possible to limit movement and delamination at the biointerface between implant and tissue as well as avoiding stress concentration that can lead to mechanical failure. Other important properties are the tensile and compressive strengths which quantify the maximum stresses a material can withstand before breaking and may be used to set stress limits that a device may be subject to within or external to the body. Depending on the application, it may be desirable for a biomaterial to have high strength so that it is resistant to failure when subjected to a load, however in other applications it may be beneficial for the material to be low strength. There is a careful balance between strength and stiffness that determines how robust to failure the biomaterial device is. Typically, as the elasticity of the biomaterial increases, the ultimate tensile strength will decrease and vice versa. One application where a high-strength material is undesired is in neural probes; if a high-strength material is used in these applications the tissue will always fail before the device does (under applied load) because the Young's Modulus of the dura mater and cerebral tissue is on the order of 500 Pa. When this happens, irreversible damage to the brain can occur, thus the biomaterial must have an elastic modulus less than or equal to brain tissue and a low tensile strength if an applied load is expected.
For implanted biomaterials that may experience temperature fluctuations, e.g., dental implants, ductility is important. The material must be ductile for a similar reason that the tensile strength cannot be too high, ductility allows the material to bend without fracture and also prevents the concentration of stresses in the tissue when the temperature changes. The material property of toughness is also important for dental implants as well as any other rigid, load-bearing implant such as a replacement hip joint. Toughness describes the material's ability to deform under applied stress without fracturing and having a high toughness allows biomaterial implants to last longer within the body, especially when subjected to large stress or cyclically loaded stresses, like the stresses applied to a hip joint during running.
For medical devices that are implanted or attached to the skin, another important property requiring consideration is the flexural rigidity, D. Flexural rigidity will determine how well the device surface can maintain conformal contact with the tissue surface, which is especially important for devices that are measuring tissue motion (strain), electrical signals (impedance), or are designed to stick to the skin without delaminating, as in epidermal electronics. Since flexural rigidity depends on the thickness of the material, h, to the third power (h3), it is very important that a biomaterial can be formed into thin layers in the previously mentioned applications where conformality is paramount.
Structure
The molecular composition of a biomaterial determines the physical and chemical properties of a biomaterial. These compositions create complex structures that allow the biomaterial to function, and therefore are necessary to define and understand in order to develop a biomaterial. biomaterials can be designed to replicate natural organisms, a process known as biomimetics. The structure of a biomaterial can be observed at different at different levels to better understand a materials properties and function.
Atomic structure
The arrangement of atoms and ions within a material is one of the most important structural properties of a biomaterial. The atomic structure of a material can be viewed at different levels, the sub atomic level, atomic or molecular level, as well as the ultra-structure created by the atoms and molecules. Intermolecular forces between the atoms and molecules that compose the material will determine its material and chemical properties.
The sub atomic level observes the electrical structure of an individual atom to define its interactions with other atoms and molecules. The molecular structure observes the arrangement of atoms within the material. Finally the ultra-structure observes the 3-D structure created from the atomic and molecular structures of the material. The solid-state of a material is characterized by the intramolecular bonds between the atoms and molecules that comprise the material. Types of intramolecular bonds include: ionic bonds, covalent bonds, and metallic bonds. These bonds will dictate the physical and chemical properties of the material, as well as determine the type of material (ceramic, metal, or polymer).
Microstructure
The microstructure of a material refers to the structure of an object, organism, or material as viewed at magnifications exceeding 25 times. It is composed of the different phases of form, size, and distribution of grains, pores, precipitates, etc. The majority of solid microstructures are crystalline, however some materials such as certain polymers will not crystallize when in the solid state.
Crystalline structure
Crystalline structure is the composition of ions, atoms, and molecules that are held together and ordered in a 3D shape. The main difference between a crystalline structure and an amorphous structure is the order of the components. Crystalline has the highest level of order possible in the material where amorphous structure consists of irregularities in the ordering pattern. One way to describe crystalline structures is through the crystal lattice, which is a three-dimensional representation of the location of a repeating factor (unit cell) in the structure denoted with lattices. There are 14 different configurations of atom arrangement in a crystalline structure, and are all represented under Bravais lattices.
Defects of crystalline structure
During the formation of a crystalline structure, different impurities, irregularities, and other defects can form. These imperfections can form through deformation of the solid, rapid cooling, or high energy radiation. Types of defects include point defects, line defects, as well as edge dislocation.
Macrostructure
Macrostructure refers to the overall geometric properties that will influence the force at failure, stiffness, bending, stress distribution, and the weight of the material. It requires little to no magnification to reveal the macrostructure of a material. Observing the macrostructure reveals properties such as cavities, porosity, gas bubbles, stratification, and fissures. The material's strength and elastic modulus are both independent of the macrostructure.
Natural biomaterials
Biomaterials can be constructed using only materials sourced from plants and animals in order to alter, replace, or repair human tissue/organs. Use of natural biomaterials were used as early as ancient Egypt, where indigenous people used animal skin as sutures. A more modern example is a hip replacement using ivory material which was first recorded in Germany 1891.
Valuable criteria for viable natural biomaterials:
Biodegradable
Biocompatible
Able to promote cell attachment and growth
Non-toxic
Examples of natural biomaterials:
Alginate
Matrigel
Fibrin
Collagen
Myocardial tissue engineering
Biopolymers
Biopolymers are polymers produced by living organisms. Cellulose and starch, proteins and peptides, and DNA and RNA are all examples of biopolymers, in which the monomeric units, respectively, are sugars, amino acids, and nucleotides. Cellulose is both the most common biopolymer and the most common organic compound on Earth. About 33% of all plant matter is cellulose. On a similar manner, silk (proteinaceous biopolymer) has garnered tremendous research interest in a myriad of domains including tissue engineering and regenerative medicine, microfluidics, drug delivery.
| Technology | Biotechnology | null |
18619339 | https://en.wikipedia.org/wiki/Slug | Slug | Slug, or land slug, is a common name for any apparently shell-less terrestrial gastropod mollusc. The word slug is also often used as part of the common name of any gastropod mollusc that has no shell, a very reduced shell, or only a small internal shell, particularly sea slugs and semi-slugs (this is in contrast to the common name snail, which applies to gastropods that have a coiled shell large enough that they can fully retract their soft parts into it).
Various taxonomic families of land slugs form part of several quite different evolutionary lineages, which also include snails. Thus, the various families of slugs are not closely related, despite the superficial similarity in overall body form. The shell-less condition has arisen many times independently as an example of convergent evolution, and thus the category "slug" is polyphyletic.
Taxonomy
Of the six orders of Pulmonata, two – the Onchidiacea and Soleolifera – solely comprise slugs. A third group, the Sigmurethra, contains various clades of snails, semi-slugs (i.e. snails whose shells are too small for them to retract fully into), and slugs. The taxonomy of this group is in the process of being revised in light of DNA sequencing. Research suggests that pulmonates are paraphyletic and basal to the opisthobranchs, which are a terminal branch of the tree. The family Ellobiidae are also polyphyletic.
Subinfraorder Orthurethra
Superfamily Achatinelloidea Gulick, 1873
Superfamily Cochlicopoidea Pilsbry, 1900
Superfamily Partuloidea Pilsbry, 1900
Superfamily Pupilloidea Turton, 1831
Subinfraorder Sigmurethra
Superfamily Acavoidea Pilsbry, 1895
Superfamily Achatinoidea Swainson, 1840
Superfamily Aillyoidea Baker, 1960
Superfamily Arionoidea J.E. Gray in Turnton, 1840
Superfamily Athoracophoroidea
Family Athoracophoridae
Superfamily Orthalicoidea
Subfamily Bulimulinae
Superfamily Camaenoidea Pilsbry, 1895
Superfamily Clausilioidea Mörch, 1864
Superfamily Dyakioidea Gude & Woodward, 1921
Superfamily Gastrodontoidea Tryon, 1866
Superfamily Helicoidea Rafinesque, 1815
Superfamily Helixarionoidea Bourguignat, 1877
Superfamily Limacoidea Rafinesque, 1815
Superfamily Oleacinoidea H. & A. Adams, 1855
Superfamily Orthalicoidea Albers-Martens, 1860
Superfamily Plectopylidoidea Moellendorf, 1900
Superfamily Polygyroidea Pilsbry, 1894
Superfamily Punctoidea Morse, 1864
Superfamily Rhytidoidea Pilsbry, 1893
Family Rhytididae
Superfamily Sagdidoidera Pilsbry, 1895
Superfamily Staffordioidea Thiele, 1931
Superfamily Streptaxoidea J.E. Gray, 1806
Superfamily Strophocheiloidea Thiele, 1926
Superfamily Parmacelloidea
Superfamily Zonitoidea Mörch, 1864
Superfamily Quijotoidea Jesús Ortea and Juan José Bacallado, 2016
Family Quijotidae
Description
The external anatomy of a slug includes the following:
Tentacles: Like other pulmonate land gastropods, the majority of land slugs have two pairs of 'feelers' or tentacles on their head. The upper pair is light-sensing and has eyespots at the ends, while the lower pair provides the sense of smell. Both pairs are retractable in stylommatophoran slugs, but contractile in veronicellid slugs.
Mantle: On top of the slug, behind the head, is the saddle-shaped mantle. In stylommatophoran slugs, on the right-hand side of the mantle is a respiratory opening, the pneumostome, which is easier to see when open; also on the right side at the front are the genital opening and anus. Veronicellid slugs have a mantle covering the whole dorsal part of the body, they have no respiratory opening, and the anus opens posteriorly.
Tail: The part of a slug behind the mantle is called the 'tail'.
Keel: Some species of slugs, for example Tandonia budapestensis, have a prominent ridge running over their back along the middle of the tail (sometimes along the whole tail, sometimes only the posterior part).
Foot: The bottom side of a slug, which is flat, is called the 'foot'. Like almost all gastropods, a slug moves by rhythmic waves of muscular contraction on the underside of its foot. It simultaneously secretes a layer of mucus that it travels on, which helps prevent damage to the foot tissues. Around the edge of the foot in some slugs is a structure called the 'foot fringe'.
Vestigial shell: Most slugs retain a remnant of their shell, which is usually internalized. This organ generally serves as storage for calcium salts, often in conjunction with the digestive glands. An internal shell is present in the Limacidae and Parmacellidae. Adult Philomycidae, Onchidiidae and Veronicellidae lack shells.
Physiology
Slugs' bodies are made up mostly of water and, without a full-sized shell, their soft tissues are prone to desiccation. They must generate protective mucus to survive. Many species are most active following rainfall or during nighttime since there is increased moisture on the ground. In drier conditions, they hide in damp places such as under tree bark, fallen logs, rocks and manmade structures, such as planters, to help retain body moisture. Like all other gastropods, they undergo torsion (a 180° twisting of the internal organs) during development. Internally, slug anatomy clearly shows the effects of this rotation—but externally, the bodies of slugs appear more or less symmetrical, except the pneumostome, which is on one side of the animal, normally the right-hand side.
Slugs produce two types of mucus: one is thin and watery, and the other thick and sticky. Both kinds are hygroscopic. The thin mucus spreads from the foot's centre to its edges, whereas the thick mucus spreads from front to back. Slugs also produce thick mucus that coats the whole body of the animal. The mucus secreted by the foot contains fibres that help prevent the slug from slipping down vertical surfaces.
The "slime trail" a slug leaves behind has some secondary effects: other slugs coming across a slime trail can recognise the slime trail as produced by one of the same species, which is useful in finding a mate. Following a slime trail is also part of the hunting behaviour of some carnivorous slugs. Body mucus provides some protection against predators, as it can make the slug hard to pick up and hold by a bird's beak, for example, or the mucus itself can be distasteful. Some slugs can also produce very sticky mucus which can incapacitate predators and can trap them within the secretion. Some species of slug, such as Limax maximus, secrete slime cords to suspend a pair during copulation.
Reproduction
Slugs are hermaphrodites, having both female and male reproductive organs. Once a slug has located a mate, they encircle each other and sperm is exchanged through their protruded genitalia.
Apophallation has been reported only in some species of banana slug (Ariolimax) and one species of Deroceras. In the banana slugs, the penis sometimes becomes trapped inside the body of the partner. Apophallation allows the slugs to separate themselves by one or both of the slugs chewing off the other's or its own penis. Once the penis has been discarded, banana slugs are still able to mate using only the female parts of the reproductive system.
In a temperate climate, slugs usually live one year outdoors. In greenhouses, many adult slugs may live for more than one year.
Ecology
Slugs play an important role in the ecosystem by eating decaying plant material and fungi. Most carnivorous slugs on occasion also eat dead specimens of their own kind.
Feeding habits
Most species of slugs are generalists, feeding on a broad spectrum of organic materials, including leaves from living plants, lichens, mushrooms, and even carrion. Some slugs are predators and eat other slugs and snails, or earthworms.
Slugs can feed on a wide variety of vegetables and herbs, including flowers such as petunias, chrysanthemums, daisies, lobelia, lilies, dahlias, narcissus, gentians, primroses, tuberous begonias, hollyhocks, marigolds, and fruits such as strawberries. They also feed on carrots, peas, apples, and cabbage that are offered as a sole food source.
Slugs from different families are fungivores. It is the case in the Philomycidae (e. g. Philomycus carolinianus and Phylomicus flexuolaris) and Ariolimacidae (Ariolimax californianus), which respectively feed on slime molds (myxomycetes) and mushrooms (basidiomycetes). Species of mushroom producing fungi used as food source by slugs include milk-caps (Lactarius spp.), the oyster mushroom (Pleurotus ostreatus) and the penny bun (Boletus edulis). Other genera such as Agaricus, Pleurocybella and Russula are also eaten by slugs. Slime molds used as food source by slugs include Stemonitis axifera and Symphytocarpus flaccidus. Some slugs are selective towards certain parts or developmental stages of the fungi they eat, though this is very variable. Depending on the species and other factors, slugs eat only fungi at specific stages of development. In other cases, whole mushrooms can be eaten, without any selection or bias towards ontogenetic stages.
Predators
Slugs are preyed upon by various vertebrates and invertebrates. The predation of slugs has been the subject of studies for at least a century. Because some species of slugs are considered agricultural pests, research investments have been made to discover and investigate potential predators in order to establish biological control strategies.
Vertebrates
Slugs are preyed upon by virtually every major vertebrate group. With many examples among reptiles, birds, mammals, amphibians and fish, vertebrates can occasionally feed on, or be specialised predators of, slugs. Fish that feed on slugs include the brown trout (Salmo trutta), which occasionally feeds on Arion circumscriptus, an arionid slug. Similarly, the shortjaw kokopu (Galaxias postvectis) includes slugs in its diet. Amphibians such as frogs and toads have long been regarded as important predators of slugs. Among them are species in the genus Bufo, Rhinella and Ceratophrys.
Reptiles that feed on slugs include mainly snakes and lizards. Some colubrid snakes are known predators of slugs. Coastal populations of the garter snake, Thamnophis elegans, have a specialised diet consisting of slugs, such as Ariolimax, while inland populations have a generalized diet. One of its congeners, the Northwestern garter snake (Thamnophis ordinoides), is not a specialized predator of slugs but occasionally feeds on them. The redbelly snake (Storeria occipitomaculata) and the brown snake (Storeria dekayi) feed mainly but not solely on slugs, while some species in the genus Dipsas (e.g. Dipsas neuwiedi) and the common slug eater snake (Duberria lutrix), are exclusively slug eaters. Several lizards include slugs in their diet. This is the case in the slowworm (Anguis fragilis), the bobtail lizard (Tiliqua rugosa), the she-oak skink (Cyclodomorphus casuarinae) and the common lizard (Zootoca vivipara).
Birds that prey upon slugs include common blackbirds (Turdus merula), starlings (Sturnus vulgaris), rooks (Corvus frugilegus), jackdaws (Corvus monedula), owls, vultures and ducks. Studies on slug predation also cite fieldfares (feeding on Deroceras reticulatum), redwings (feeding on Limax and Arion), thrushes (on Limax and Arion ater), red grouse (on Deroceras and Arion hortensis), game birds, wrynecks (on Limax flavus), rock doves and charadriiform birds as slug predators.
Mammals that eat slugs include foxes, badgers and hedgehogs.
Invertebrates
Beetles in the family Carabidae, such as Carabus violaceus and Pterostichus melanarius, are known to feed on slugs. Ants are a common predator of slugs; some ant species are deterred by the slug's mucus coating, while others such as driver ants will roll the slug in dirt to absorb its mucus.
Parasites and parasitoids
Slugs are parasitised by several organisms, including acari and a wide variety of nematodes. The slug mite, Riccardoella limacum, is known to parasitise several dozen species of molluscs, including many slugs, such as Deroceras reticulatum, Arianta arbustorum, Arion ater, Arion hortensis, Limax maximus, Tandonia budapestensis, Milax gagates, and Tandonia sowerbyi. R. limacum can often be seen swarming about their host's body, and live in its respiratory cavity.
Several species of nematodes are known to parasitise slugs. The nematode worms Agfa flexilis and Angiostoma limacis respectively live in the salivary glands and rectum of Limax maximus. Species of widely known medical importance pertaining to the genus Angiostrongylus are also parasites of slugs. Both Angiostrongylus costaricensis and Angiostrongylus cantonensis, a meningitis-causing nematode, have larval stages that can only live in molluscs, including slugs, such as Limax maximus.
Insects such as dipterans are known parasitoids of molluscs. To complete their development, many dipterans use slugs as hosts during their ontogeny. Some species of blow-flies (Calliphoridae) in the genus Melinda are known parasitoids of Arionidae, Limacidae and Philomycidae. Flies in the family Phoridae, specially those in the genus Megaselia, are parasitoids of Agriolimacidae, including many species of Deroceras. House flies in the family Muscidae, mainly those in the genus Sarcophaga, are facultative parasitoids of Arionidae.
Behavior
When attacked, slugs can contract their body, making themselves harder and more compact and more still and round. By doing this, they become firmly attached to the substrate. This, combined with the slippery mucus they produce, makes slugs more difficult for predators to grasp. The unpleasant taste of the mucus is also a deterrent. Slugs can also incapacitate predators through the production of a highly sticky and elastic mucus which can trap predators in the secretion.
Some species present different response behaviors when attacked, such as the Kerry slug. In contrast to the general behavioral pattern, the Kerry slug retracts its head, lets go of the substrate, rolls up completely, and stays contracted in a ball-like shape. This is a unique feature among all the Arionidae, and among most other slugs. Some slugs can self-amputate (autotomy) a portion of their tail to help the slug escape from a predator. Some slug species hibernate underground during the winter in temperate climates, but in other species, the adults die in the autumn.
Intra- and inter-specific agonistic behavior is documented, but varies greatly among slug species. Slugs often resort to aggression, attacking both conspecifics and individuals from other species when competing for resources. This aggressiveness is also influenced by seasonality, because the availability of resources such as shelter and food may be compromised due to climatic conditions. Slugs are prone to attack during the summer, when the availability of resources is reduced. During winter, the aggressive responses are substituted by a gregarious behavior.
Human relevance
The great majority of slug species are harmless to humans and to their interests, but a small number of species are serious pests of agriculture and horticulture. They can destroy foliage faster than plants can grow, thus killing even fairly large plants. They also feed on fruits and vegetables prior to harvest, making holes in the crop, which can make individual items unsuitable to sell for aesthetic reasons, and can make the crop more vulnerable to rot and disease. Excessive buildup of slugs within some wastewater treatment plants with inadequate screening have been found to cause process issues resulting in increased energy and chemical use.
In a few rare cases, humans have developed Angiostrongylus cantonensis-induced meningitis from eating raw slugs. Live slugs that are accidentally eaten with improperly cleaned vegetables (such as lettuce), or improperly cooked slugs (for use in recipes requiring larger slugs such as banana slugs), can act as a vector for a parasitic infection in humans.
Prevention
As control measures, baits are commonly used in both agriculture and the garden. In recent years, iron phosphate baits have emerged and are preferred over the more toxic metaldehyde, especially because domestic or wild animals may be exposed to the bait. The environmentally safer iron phosphate has been shown to be at least as effective as baits. Methiocarb baits are no longer widely used. Parasitic nematodes (Phasmarhabditis hermaphrodita) are a commercially available biological control method that are effective against a wide range of common slug species. The nematodes are applied in water and actively seek out slugs in the soil and infect them, leading to the death of the slug. This control method is suitable for use in organic growing systems.
Other slug control methods are generally ineffective on a large scale, but can be somewhat useful in small gardens. These include , diatomaceous earth, crushed eggshells, coffee grounds, and copper. Salt kills slugs by causing water to leave the body owing to osmosis but this is not used for agricultural control as soil salinity is detrimental to crops. Conservation tillage worsens slug infestations. Hammond et al. 1999 find maize/corn and soybean in the US to be more severely affected under low till because this increases organic matter, thus providing food and shelter.
Gallery
| Biology and health sciences | Mollusks | null |
18625077 | https://en.wikipedia.org/wiki/Goldfish | Goldfish | The goldfish (Carassius auratus) is a freshwater fish in the family Cyprinidae of order Cypriniformes. It is commonly kept as a pet in indoor aquariums, and is one of the most popular aquarium fish. Goldfish released into the wild have become an invasive pest in parts of North America and Australia.
Native to China, the goldfish is a relatively small member of the carp family (which also includes the Prussian carp and the crucian carp). It was first selectively bred for color in imperial China more than 1,000 years ago, where several distinct breeds were developed. Goldfish breeds vary greatly in size, body shape, fin configuration, and coloration (various combinations of white, yellow, orange, red, brown, and black are known).
History
Various species of carp (collectively known as Asian carp) have been bred and reared as food fish for thousands of years in East Asia. Some of these normally gray or silver species have a tendency to produce red, orange, or yellow color mutations; this was first recorded in Imperial China, during the Jin dynasty (266–420).
During the Tang dynasty (AD 618–907), it was popular to raise carp in ornamental ponds and water gardens. A natural genetic mutation produced gold (actually yellowish orange) rather than silver coloration. People began to selectively breed the gold variety instead of the silver variety, keeping them in ponds or other bodies of water. On special occasions at which guests were expected, they would be moved to a much smaller container for display.
By the Song dynasty (AD 960–1279), the selective domestic breeding of goldfish was firmly established. In 1162, the empress of the Song dynasty ordered the construction of a pond to collect the red and gold variety. By this time, people outside the imperial family were forbidden to keep goldfish of the gold (yellow) variety, yellow being the imperial color.
During the Ming dynasty (1368–1644), goldfish also began to be raised indoors, which permitted selection for mutations that would not be able to survive in ponds. The first occurrence of -tailed goldfish was recorded in the Ming dynasty. In 1603, goldfish were introduced to Japan. In 1611, goldfish were introduced to Portugal and from there to other parts of Europe.
During the 1620s, goldfish were highly regarded in southern Europe because of their metallic scales, and symbolized good luck and fortune. It became a tradition for married men to give their wives a goldfish on their first anniversary, as a symbol for the prosperous years to come. This tradition quickly died, as goldfish became more available, losing their status. Goldfish were first introduced to North America around 1850 and quickly became popular in the United States.
Biology
Taxonomy
There has been considerable debate about the taxonomy of the goldfish. Previously, the goldfish was believed to be either a subspecies of the crucian carp (Carassius carassius), or of the Prussian carp (Carassius gibelio). However, modern genetic sequencing has suggested otherwise, and that modern goldfish are domesticated varieties of C. auratus that are native to Southern China. C. auratus are differentiated from other Carassius species by several characteristics. C. auratus have a more pointed snout, while the snout of C. carassius is well rounded. C. gibelio often has a grayish/greenish color, while crucian carp are always golden bronze. Juvenile crucian carp have a black spot on the base of the tail, which disappears with age. In C. auratus, this tail spot is never present. C. auratus have fewer than 31 scales along the lateral line, while crucian carp have 33 scales or more.
Goldfish can hybridize with some other Carassius species of carp. Koi and common carp may also interbreed with goldfish to produce sterile hybrids.
Size
Wild goldfish typically grow to between 4.7 inches (12 cm) and 8.7 inches (22 cm) but can reach 16 inches (41 cm). The size of pet goldfish depends upon its breed.
As of April 2008, the largest goldfish in the world was believed by the BBC to measure , in the Netherlands. At the time, a goldfish named "Goldie", kept as a pet in a tank in Folkestone, England, was measured as and over , and named as the second largest in the world behind the Netherlands fish. The secretary of the Federation of British Aquatic Societies (FBAS) stated of Goldie's size, "I would think there are probably a few bigger goldfish that people don't think of as record holders, perhaps in ornamental lakes". In July 2010, a goldfish measuring and was caught in a pond in Poole, England, thought to have been abandoned there after outgrowing a tank. On November 16, 2020, a goldfish weighing was found in a lake in Greenville, South Carolina, while conducting a population survey of Oak Grove Lake.
Vision
As a domestic fish, thus an easily accessible model organism, goldfish have one of the most studied senses of vision in fishes. Goldfish have four kinds of cone cells, which are respectively sensitive to different colors: red, green, blue and ultraviolet. The ability to distinguish between four different primary colors classifies them as tetrachromats.
Hearing
Goldfish have one of the most studied senses of hearing in fish. They have two otoliths, permitting the detection of sound particle motion, and Weberian ossicles connecting the swim bladder to the otoliths, facilitating the detection of sound pressure.
Reproduction
Goldfish can only grow to sexual maturity with enough water and the right nutrition. Most goldfish breed in captivity, particularly in pond settings. Breeding usually happens after a significant temperature change, often in spring. Males chase gravid female goldfish (females carrying eggs), and prompt them to release their eggs by bumping and nudging them.
Goldfish, like all cyprinids, are egg-layers. Their eggs are adhesive and attach to aquatic vegetation, typically dense plants such as Cabomba or Elodea or a spawning mop. The eggs hatch within 48 to 72 hours.
Within a week or so, the fry begins to assume its final shape, although a year may pass before they develop a mature goldfish color; until then they are a metallic brown like their wild ancestors. In their first weeks of life, the fry grow quickly—an adaptation born of the high risk of getting devoured by the adult goldfish (or other fish and insects) in their environment.
Some highly selectively bred goldfish can no longer breed naturally due to their altered shape. The artificial breeding method called "hand stripping" can assist in breeding, but can harm the fish if not done correctly. In captivity, adults may also eat young that they encounter.
Respiration
Goldfish are able to survive short periods of entirely anoxic conditions. Survival is shorter under higher temperatures, suggesting that this is a cold weather adaptation. Researchers speculate that this is specifically an adaptation to survival in frozen water bodies over winter.
Energy is obtained from liver glycogen. This process depends upon a pyruvate decarboxylase – the first known in vertebrates.
Salinity
Although they are a freshwater fish, goldfish have been found in brackish water with a salinity of 17.
Behavior
Goldfish are gregarious, displaying schooling behavior, as well as displaying the same types of feeding behaviors.
Goldfish have learned behaviors, both as groups and as individuals, that stem from native carp behavior. They are a generalist species with varied feeding, breeding, and predator avoidance behaviors that contribute to their success. As fish, they can be described as "friendly" towards each other. Very rarely does a goldfish harm another goldfish, nor do the males harm the females during breeding. The only real threat that goldfish present to each other is competing for food. Commons, comets, and other faster varieties can easily eat all the food during a feeding before varieties can reach it. This can lead to stunted growth or possible starvation of fancier varieties when they are kept in a pond with their single-tailed brethren. As a result, care should be taken to combine only breeds with similar body type and swim characteristics.
Cognitive abilities
Goldfish have strong associative learning abilities, as well as social learning skills. In addition, their visual acuity allows them to distinguish between individual humans. Owners may notice that fish react favorably to them (swimming to the front of the glass, swimming rapidly around the tank, and going to the surface mouthing for food) while hiding when other people approach the tank. Over time, goldfish learn to associate their owners and other humans with food, often "begging" for food whenever their owners approach.
Goldfish that have constant visual contact with humans also stop considering them to be a threat. After being kept in a tank for several weeks, sometimes months, it becomes possible to feed a goldfish by hand without it shying away.
Goldfish have a memory-span of at least three months and can distinguish between different shapes, colors, and sounds. By using positive reinforcement, goldfish can be trained to recognize and to react to light signals of different colors or to perform tricks. Fish respond to certain colors most evidently in relation to feeding. Fish learn to anticipate feedings provided they occur at around the same time every day.
Classification
Western
As with many other examples of animal, selective breeding of goldfish over centuries has produced several color variations, some of them far removed from the "golden" color of the original fish. There are also different body shapes, and fin and eye configurations. Some extreme versions of the goldfish live only in aquariums—they are much less hardy than varieties closer to the "wild" original. However, some variations are hardier, such as the Shubunkin. Currently, there are about 300 breeds recognized in China. The vast majority of goldfish breeds today originated from China. Some of the main varieties are:
Chinese
Chinese tradition classifies goldfish into four main types. These classifications are not commonly used in the West.
Crucian (also called "grass") — Goldfish without anatomical features, similar to Crucian carp or grass carp except for their coloration. These include the common goldfish, comet goldfish and Shubunkin.
Wen — Goldfish having a tail, e.g., fantails and veiltails. "Wen" is also the name of the characteristic headgrowth on such strains as oranda and lionhead.
Dragon Eye — Goldfish having extended eyes, e.g., black moor, bubble eye, and telescope eye
Egg — Goldfish having no dorsal fin, usually with an 'egg-shaped' body, e.g., lionhead. This group includes a bubble eye without a dorsal fin.
Cultivation
In aquaria
Like most species in the carp family, goldfish produce a large amount of waste both in their feces and through their gills, releasing harmful chemicals into the water. Buildup of this waste to toxic levels can occur in a relatively short period of time, and can easily cause a goldfish's death. For common and comet varieties, each goldfish should have about of water. Smaller fantail goldfish should have about per goldfish. The water surface area determines how much oxygen diffuses and dissolves into the water. A general rule is have . Active aeration by way of a water pump, filter or fountain effectively increases the surface area agitation.
The goldfish is classified as a coldwater fish, and can live in unheated aquaria at a temperature comfortable for humans. However, rapid changes in temperature, for example in an office building in winter when the heat is turned off at night, can kill them, especially if the tank is small. Care must also be taken when adding water, as the new water may be of a different temperature. Temperatures under about are dangerous to fancy varieties, though commons and comets can survive slightly lower temperatures. Extremely high temperatures (over ) can also harm goldfish. However, higher temperatures may help fight protozoan infestations by accelerating the parasite's life cycle—thus eliminating it more quickly. The optimum temperature for goldfish is between .
Like all fish, goldfish do not like to be petted. In fact, touching a goldfish can endanger its health, because it can cause the protective slime coat to be damaged or removed, exposing the fish's skin to infection from bacteria or water-borne parasites. However, goldfish respond to people by surfacing at feeding time, and can be trained or acclimated to taking pellets or flakes from human fingers. The reputation of goldfish dying quickly is often due to poor care. The lifespan of goldfish in captivity can extend beyond 10 years.
If left in the dark for a period of time, goldfish gradually change color until they are almost gray. Goldfish produce pigment in response to light, similarly to how human skin becomes tanned in the sun. Fish have cells called chromatophores that produce pigments that reflect light and give the fish coloration. The color of a goldfish is determined by their diet, water quality, and exposure to light, along with age and health.
Because goldfish eat live plants, their presence in a planted aquarium can be problematic. Only a few aquarium plant species, such as Cryptocoryne and Anubias, can survive around goldfish, but they require special attention so that they are not uprooted.
In ponds
Goldfish are popular pond fish, since they are small, inexpensive, colorful, and very hardy. In an outdoor pond or water garden, they may even survive for brief periods if ice forms on the surface, as long as there is enough oxygen remaining in the water and the pond does not freeze solid. Common, London and Bristol shubunkins, jikin, wakin, comet and some hardier fantail goldfish can be kept in a pond year-round in temperate and subtropical climates. Moor, veiltail, oranda and lionhead can be kept safely in outdoor ponds year-round only in more tropical climates and elsewhere only in summer months.
Compatible fish include rudd, tench, orfe and koi, but the last require specialized care. Ramshorn snails are helpful by eating any algae that grows in the pond. Without some form of animal population control, goldfish ponds can easily become overstocked. Fish such as orfe consume goldfish eggs.
Ponds small and large are fine in warmer areas, though goldfish can "overheat" in small volumes of water in the summer in tropical climates. In frosty climes, the depth should be at least to preclude freezing. During winter, goldfish become sluggish, stop eating and often stay on the bottom of the pond. This is normal; they become active again in the spring. Unless the pond is large enough to maintain its own ecosystem without interference from humans, a filter is important to clear waste and keep the pond clean. Plants are essential as they act as part of the filtration system, as well as a food source for the fish. Plants are further beneficial since they raise oxygen levels in the water.
Like their wild ancestors, common and comet goldfish as well as shubunkin can survive, and even thrive, in any climate that can support a pond. In general, when released into the wild, goldfish quickly take over the waterways as an invasive species.
Feeding
In the wild, the diet of goldfish consists of crustaceans, insects, and various plant matter. Like most fish, they are opportunistic feeders and do not stop eating on their own accord. Overfeeding can be deleterious to their health, typically by blocking the intestines. This happens most often with selectively bred goldfish, which have a convoluted intestinal tract. When excess food is available, they produce more waste and feces, partly due to incomplete protein digestion. Overfeeding can sometimes be diagnosed by observing feces trailing from the fish's cloaca.
Goldfish-specific food has less protein and more carbohydrate than conventional fish food. Enthusiasts may supplement this diet with shelled peas (with outer skins removed), blanched green leafy vegetables, and bloodworms. Young goldfish benefit from the addition of brine shrimp to their diet. As with all animals, goldfish preferences vary.
For mosquito control
Like some other well-known aquarium fish, such as the guppy and mosquitofish, goldfish (and other carp) are frequently added to stagnant bodies of water in an attempt to reduce mosquito populations, which spread the vectors of diseases such as West Nile virus, malaria, and dengue. However, introducing goldfish has often had negative consequences for local ecosystems, and their efficacy as pest control has never been compared to those of native fishes.
Market
The market for live goldfish and other crucian carp usually imported from China was $1.2million in 2018. Some high quality varieties cost between $125 and $300.
Welfare concerns
Fishbowls are detrimental to the health of goldfish and are prohibited by animal welfare legislation in several municipalities. The practice of using bowls as permanent fish housing originated from a misunderstanding of Chinese "display" vessels: goldfish which were normally housed in ponds were, on occasion, temporarily displayed in smaller containers to be better admired by guests.
Goldfish kept in bowls or "mini-aquariums" suffer from death, disease, and stunting, due primarily to the low oxygen and very high ammonia/nitrite levels inherent in such an environment. In comparison to other common aquarium fish, goldfish have high oxygen needs and produce a large amount of waste due to the fact they lack a stomach; therefore they require a substantial volume of well-filtered water to thrive. In addition, all goldfish varieties have the potential to reach in total length, with single-tailed breeds often exceeding . Single-tailed varieties include common and comet goldfish.
In many countries, carnival and fair operators commonly give goldfish away in plastic bags as prizes. In late 2005 Rome banned the use of goldfish and other animals as carnival prizes. Rome has also banned the use of "goldfish bowls", on animal cruelty grounds, as well as Monza, Italy, in 2004. In the United Kingdom, the government proposed banning this practice as part of its Animal Welfare Bill, though this has since been amended to only prevent goldfish being given as prizes to unaccompanied minors.
In Japan, during summer festivals and religious holidays (ennichi), a traditional game called goldfish scooping is played, in which a player scoops goldfish from a basin with a special scooper. Sometimes bouncy balls are substituted for goldfish.
Although edible and closely related to some fairly widely eaten species, goldfish are rarely eaten. A fad among American college students for many years was swallowing goldfish as a stunt and as a fraternity initiation process. The first recorded instance was in 1939 at Harvard University. The practice gradually fell out of popularity over the course of several decades and is rarely practiced today.
Some animal advocates have called for boycotts of goldfish purchases, citing industrial farming and low survival rates of the fish.
In popular culture
In Chinese history, goldfish was seen "as a symbol of luck and fortune". Moreover, only members of the Song dynasty could own goldfish. In Iran and among the international Iranian diaspora, goldfish is a traditional part of Nowruz celebrations. Goldfish are usually placed on Haft-sin tables as a symbol of progress.
| Biology and health sciences | Cypriniformes | null |
23868856 | https://en.wikipedia.org/wiki/Reynolds%20number | Reynolds number | In fluid dynamics, the Reynolds number () is a dimensionless quantity that helps predict fluid flow patterns in different situations by measuring the ratio between inertial and viscous forces. At low Reynolds numbers, flows tend to be dominated by laminar (sheet-like) flow, while at high Reynolds numbers, flows tend to be turbulent. The turbulence results from differences in the fluid's speed and direction, which may sometimes intersect or even move counter to the overall direction of the flow (eddy currents). These eddy currents begin to churn the flow, using up energy in the process, which for liquids increases the chances of cavitation.
The Reynolds number has wide applications, ranging from liquid flow in a pipe to the passage of air over an aircraft wing. It is used to predict the transition from laminar to turbulent flow and is used in the scaling of similar but different-sized flow situations, such as between an aircraft model in a wind tunnel and the full-size version. The predictions of the onset of turbulence and the ability to calculate scaling effects can be used to help predict fluid behavior on a larger scale, such as in local or global air or water movement, and thereby the associated meteorological and climatological effects.
The concept was introduced by George Stokes in 1851, but the Reynolds number was named by Arnold Sommerfeld in 1908 after Osborne Reynolds who popularized its use in 1883 (an example of Stigler's law of eponymy).
Definition
The Reynolds number is the ratio of inertial forces to viscous forces within a fluid that is subjected to relative internal movement due to different fluid velocities. A region where these forces change behavior is known as a boundary layer, such as the bounding surface in the interior of a pipe. A similar effect is created by the introduction of a stream of high-velocity fluid into a low-velocity fluid, such as the hot gases emitted from a flame in air. This relative movement generates fluid friction, which is a factor in developing turbulent flow. Counteracting this effect is the viscosity of the fluid, which tends to inhibit turbulence. The Reynolds number quantifies the relative importance of these two types of forces for given flow conditions and is a guide to when turbulent flow will occur in a particular situation.
This ability to predict the onset of turbulent flow is an important design tool for equipment such as piping systems or aircraft wings, but the Reynolds number is also used in scaling of fluid dynamics problems and is used to determine dynamic similitude between two different cases of fluid flow, such as between a model aircraft, and its full-size version. Such scaling is not linear and the application of Reynolds numbers to both situations allows scaling factors to be developed.
With respect to laminar and turbulent flow regimes:
laminar flow occurs at low Reynolds numbers, where viscous forces are dominant, and is characterized by smooth, constant fluid motion;
turbulent flow occurs at high Reynolds numbers and is dominated by inertial forces, which tend to produce chaotic eddies, vortices and other flow instabilities.
The Reynolds number is defined as:where:
is the density of the fluid (SI units: kg/m3)
is the flow speed (m/s)
is a characteristic length (m)
is the dynamic viscosity of the fluid (Pa·s or N·s/m2 or kg/(m·s))
is the kinematic viscosity of the fluid (m2/s).
The Reynolds number can be defined for several different situations where a fluid is in relative motion to a surface. These definitions generally include the fluid properties of density and viscosity, plus a velocity and a characteristic length or characteristic dimension (L in the above equation). This dimension is a matter of convention—for example radius and diameter are equally valid to describe spheres or circles, but one is chosen by convention. For aircraft or ships, the length or width can be used. For flow in a pipe, or for a sphere moving in a fluid, the internal diameter is generally used today. Other shapes such as rectangular pipes or non-spherical objects have an equivalent diameter defined. For fluids of variable density such as compressible gases or fluids of variable viscosity such as non-Newtonian fluids, special rules apply. The velocity may also be a matter of convention in some circumstances, notably stirred vessels.
In practice, matching the Reynolds number is not on its own sufficient to guarantee similitude. Fluid flow is generally chaotic, and very small changes to shape and surface roughness of bounding surfaces can result in very different flows. Nevertheless, Reynolds numbers are a very important guide and are widely used.
Derivation
If we know that the relevant physical quantities in a physical system are only , then the Reynolds number is essentially fixed by the Buckingham π theorem.
In detail, since there are 4 quantities , but they have only 3 dimensions (length, time, mass), we can consider , where are real numbers. Setting the three dimensions of to zero, we obtain 3 independent linear constraints, so the solution space has 1 dimension, and it is spanned by the vector .
Thus, any dimensionless quantity constructed out of is a function of , the Reynolds number.
Alternatively, we can take the incompressible Navier–Stokes equations (convective form):Remove the gravity term , then the left side consists of inertial force , and viscous force . Their ratio has the order of , the Reynolds number. This argument is written out in detail on the Scallop theorem page.
Alternative derivation
The Reynolds number can be obtained when one uses the nondimensional form of the incompressible Navier–Stokes equations for a newtonian fluid expressed in terms of the Lagrangian derivative:
Each term in the above equation has the units of a "body force" (force per unit volume) with the same dimensions of a density times an acceleration. Each term is thus dependent on the exact measurements of a flow. When one renders the equation nondimensional, that is when we multiply it by a factor with inverse units of the base equation, we obtain a form that does not depend directly on the physical sizes. One possible way to obtain a nondimensional equation is to multiply the whole equation by the factor
where
is the mean velocity, or , relative to the fluid (m/s),
is the characteristic length (m),
is the fluid density (kg/m3).
If we now set
we can rewrite the Navier–Stokes equation without dimensions:
where the term .
Finally, dropping the primes for ease of reading:
This is why mathematically all Newtonian, incompressible flows with the same Reynolds number are comparable. Notice also that in the above equation, the viscous terms vanish for . Thus flows with high Reynolds numbers are approximately inviscid in the free stream.
History
Osborne Reynolds famously studied the conditions in which the flow of fluid in pipes transitioned from laminar flow to turbulent flow.
In his 1883 paper Reynolds described the transition from laminar to turbulent flow in a classic experiment in which he examined the behaviour of water flow under different flow velocities using a small stream of dyed water introduced into the centre of clear water flow in a larger pipe.
The larger pipe was glass so the behaviour of the layer of the dyed stream could be observed. At the end of this pipe, there was a flow control valve used to vary the water velocity inside the tube. When the velocity was low, the dyed layer remained distinct throughout the entire length of the large tube. When the velocity was increased, the layer broke up at a given point and diffused throughout the fluid's cross-section. The point at which this happened was the transition point from laminar to turbulent flow.
From these experiments came the dimensionless Reynolds number for dynamic similarity—the ratio of inertial forces to viscous forces. Reynolds also proposed what is now known as the Reynolds averaging of turbulent flows, where quantities such as velocity are expressed as the sum of mean and fluctuating components. Such averaging allows for 'bulk' description of turbulent flow, for example using the Reynolds-averaged Navier–Stokes equations.
Flow in a pipe
For flow in a pipe or tube, the Reynolds number is generally defined as
where
is the hydraulic diameter of the pipe (the inside diameter if the pipe is circular) (m),
is the volumetric flow rate (m3/s),
is the pipe's cross-sectional area () (m2),
is the mean velocity of the fluid (m/s),
(mu) is the dynamic viscosity of the fluid (Pa·s = N·s/m2 = kg/(m·s)),
(nu) is the kinematic viscosity () (m2/s),
(rho) is the density of the fluid (kg/m3),
is the mass flowrate of the fluid (kg/s).
For shapes such as squares, rectangular or annular ducts where the height and width are comparable, the characteristic dimension for internal-flow situations is taken to be the hydraulic diameter, , defined as
where is the cross-sectional area, and is the wetted perimeter. The wetted perimeter for a channel is the total perimeter of all channel walls that are in contact with the flow. This means that the length of the channel exposed to air is not included in the wetted perimeter.
For a circular pipe, the hydraulic diameter is exactly equal to the inside pipe diameter:
For an annular duct, such as the outer channel in a tube-in-tube heat exchanger, the hydraulic diameter can be shown algebraically to reduce to
where
is the inside diameter of the outer pipe,
is the outside diameter of the inner pipe.
For calculation involving flow in non-circular ducts, the hydraulic diameter can be substituted for the diameter of a circular duct, with reasonable accuracy, if the aspect ratio AR of the duct cross-section remains in the range < AR < 4.
Laminar–turbulent transition
In boundary layer flow over a flat plate, experiments confirm that, after a certain length of flow, a laminar boundary layer will become unstable and turbulent. This instability occurs across different scales and with different fluids, usually when ≈ , where is the distance from the leading edge of the flat plate, and the flow velocity is the freestream velocity of the fluid outside the boundary layer.
For flow in a pipe of diameter , experimental observations show that for "fully developed" flow, laminar flow occurs when < 2300 and turbulent flow occurs when > 2900. At the lower end of this range, a continuous turbulent-flow will form, but only at a very long distance from the inlet of the pipe. The flow in between will begin to transition from laminar to turbulent and then back to laminar at irregular intervals, called intermittent flow. This is due to the different speeds and conditions of the fluid in different areas of the pipe's cross-section, depending on other factors such as pipe roughness and flow uniformity. Laminar flow tends to dominate in the fast-moving center of the pipe while slower-moving turbulent flow dominates near the wall. As the Reynolds number increases, the continuous turbulent-flow moves closer to the inlet and the intermittency in between increases, until the flow becomes fully turbulent at > 2900. This result is generalized to non-circular channels using the hydraulic diameter, allowing a transition Reynolds number to be calculated for other shapes of channel.
These transition Reynolds numbers are also called critical Reynolds numbers, and were studied by Osborne Reynolds around 1895. The critical Reynolds number is different for every geometry.
Flow in a wide duct
For a fluid moving between two plane parallel surfaces—where the width is much greater than the space between the plates—then the characteristic dimension is equal to the distance between the plates. This is consistent with the annular duct and rectangular duct cases above, taken to a limiting aspect ratio.
Flow in an open channel
For calculating the flow of liquid with a free surface, the hydraulic radius must be determined. This is the cross-sectional area of the channel divided by the wetted perimeter. For a semi-circular channel, it is a quarter of the diameter (in case of full pipe flow). For a rectangular channel, the hydraulic radius is the cross-sectional area divided by the wetted perimeter. Some texts then use a characteristic dimension that is four times the hydraulic radius, chosen because it gives the same value of for the onset of turbulence as in pipe flow, while others use the hydraulic radius as the characteristic length-scale with consequently different values of for transition and turbulent flow.
Flow around airfoils
Reynolds numbers are used in airfoil design to (among other things) manage "scale effect" when computing/comparing characteristics (a tiny wing, scaled to be huge, will perform differently). Fluid dynamicists define the chord Reynolds number , where is the flight speed, is the chord length, and is the kinematic viscosity of the fluid in which the airfoil operates, which is for the atmosphere at sea level. In some special studies a characteristic length other than chord may be used; rare is the "span Reynolds number", which is not to be confused with span-wise stations on a wing, where chord is still used.
Object in a fluid
The Reynolds number for an object moving in a fluid, called the particle Reynolds number and often denoted , characterizes the nature of the surrounding flow and its fall velocity.
In viscous fluids
Where the viscosity is naturally high, such as polymer solutions and polymer melts, flow is normally laminar. The Reynolds number is very small and Stokes' law can be used to measure the viscosity of the fluid. Spheres are allowed to fall through the fluid and they reach the terminal velocity quickly, from which the viscosity can be determined.
The laminar flow of polymer solutions is exploited by animals such as fish and dolphins, who exude viscous solutions from their skin to aid flow over their bodies while swimming. It has been used in yacht racing by owners who want to gain a speed advantage by pumping a polymer solution such as low molecular weight polyoxyethylene in water, over the wetted surface of the hull.
It is, however, a problem for mixing polymers, because turbulence is needed to distribute fine filler (for example) through the material. Inventions such as the "cavity transfer mixer" have been developed to produce multiple folds into a moving melt so as to improve mixing efficiency. The device can be fitted onto extruders to aid mixing.
Sphere in a fluid
For a sphere in a fluid, the characteristic length-scale is the diameter of the sphere and the characteristic velocity is that of the sphere relative to the fluid some distance away from the sphere, such that the motion of the sphere does not disturb that reference parcel of fluid. The density and viscosity are those belonging to the fluid. Note that purely laminar flow only exists up to = 10 under this definition.
Under the condition of low , the relationship between force and speed of motion is given by Stokes' law.
At higher Reynolds numbers the drag on a sphere depends on surface roughness. Thus, for example, adding dimples on the surface of a golf ball causes the boundary layer on the upstream side of the ball to transition from laminar to turbulent. The turbulent boundary layer is able to remain attached to the surface of the ball much longer than a laminar boundary and so creates a narrower low-pressure wake and hence less pressure drag. The reduction in pressure drag causes the ball to travel farther.
Rectangular object in a fluid
The equation for a rectangular object is identical to that of a sphere, with the object being approximated as an ellipsoid and the axis of length being chosen as the characteristic length scale. Such considerations are important in natural streams, for example, where there are few perfectly spherical grains. For grains in which measurement of each axis is impractical, sieve diameters are used instead as the characteristic particle length-scale. Both approximations alter the values of the critical Reynolds number.
Fall velocity
The particle Reynolds number is important in determining the fall velocity of a particle. When the particle Reynolds number indicates laminar flow, Stokes' law can be used to calculate its fall velocity or settling velocity. When the particle Reynolds number indicates turbulent flow, a turbulent drag law must be constructed to model the appropriate settling velocity.
Packed bed
For fluid flow through a bed, of approximately spherical particles of diameter in contact, if the voidage is and the superficial velocity is , the Reynolds number can be defined as
or
or
The choice of equation depends on the system involved: the first is successful in correlating the data for various types of packed and fluidized beds, the second Reynolds number suits for the liquid-phase data, while the third was found successful in correlating the fluidized bed data, being first introduced for liquid fluidized bed system.
Laminar conditions apply up to = 10, fully turbulent from = 2000.
Stirred vessel
In a cylindrical vessel stirred by a central rotating paddle, turbine or propeller, the characteristic dimension is the diameter of the agitator . The velocity is where is the rotational speed in rad per second. Then the Reynolds number is:
The system is fully turbulent for values of above .
Pipe friction
Pressure drops seen for fully developed flow of fluids through pipes can be predicted using the Moody diagram which plots the Darcy–Weisbach friction factor against Reynolds number and relative roughness . The diagram clearly shows the laminar, transition, and turbulent flow regimes as Reynolds number increases. The nature of pipe flow is strongly dependent on whether the flow is laminar or turbulent.
Similarity of flows
In order for two flows to be similar, they must have the same geometry and equal Reynolds and Euler numbers. When comparing fluid behavior at corresponding points in a model and a full-scale flow, the following holds:
where is the Reynolds number for the model, and is full-scale Reynolds number, and similarly for the Euler numbers.
The model numbers and design numbers should be in the same proportion, hence
This allows engineers to perform experiments with reduced scale models in water channels or wind tunnels and correlate the data to the actual flows, saving on costs during experimentation and on lab time. Note that true dynamic similitude may require matching other dimensionless numbers as well, such as the Mach number used in compressible flows, or the Froude number that governs open-channel flows. Some flows involve more dimensionless parameters than can be practically satisfied with the available apparatus and fluids, so one is forced to decide which parameters are most important. For experimental flow modeling to be useful, it requires a fair amount of experience and judgment of the engineer.
An example where the mere Reynolds number is not sufficient for the similarity of flows (or even the flow regime – laminar or turbulent) are bounded flows, i.e. flows that are restricted by walls or other boundaries. A classical example of this is the Taylor–Couette flow, where the dimensionless ratio of radii of bounding cylinders is also important, and many technical applications where these distinctions play an important role. Principles of these restrictions were developed by Maurice Marie Alfred Couette and Geoffrey Ingram Taylor and developed further by Floris Takens and David Ruelle.
Typical values of Reynolds number
Dictyostelium amoebae: ~ 1 × 10−6
Bacterium ~ 1 × 10−4
Ciliate ~ 1 × 10−1
Smallest fish ~ 1
Blood flow in brain ~ 1 × 102
Blood flow in aorta ~ 1 × 103
Onset of turbulent flow ~ 2.3 × 103 to 5.0 × 104 for pipe flow to 106 for boundary layers
Typical pitch in Major League Baseball ~ 2 × 105
Person swimming ~ 4 × 106
Fastest fish ~ 1 × 108
Blue whale ~ 4 × 108
A large ship (Queen Elizabeth 2) ~ 5 × 109
Atmospheric tropical cyclone ~ 1 x 1012
Smallest scales of turbulent motion
In a turbulent flow, there is a range of scales of the time-varying fluid motion. The size of the largest scales of fluid motion (sometimes called eddies) are set by the overall geometry of the flow. For instance, in an industrial smoke stack, the largest scales of fluid motion are as big as the diameter of the stack itself. The size of the smallest scales is set by the Reynolds number. As the Reynolds number increases, smaller and smaller scales of the flow are visible. In a smokestack, the smoke may appear to have many very small velocity perturbations or eddies, in addition to large bulky eddies. In this sense, the Reynolds number is an indicator of the range of scales in the flow. The higher the Reynolds number, the greater the range of scales. The largest eddies will always be the same size; the smallest eddies are determined by the Reynolds number.
What is the explanation for this phenomenon? A large Reynolds number indicates that viscous forces are not important at large scales of the flow. With a strong predominance of inertial forces over viscous forces, the largest scales of fluid motion are undamped—there is not enough viscosity to dissipate their motions. The kinetic energy must "cascade" from these large scales to progressively smaller scales until a level is reached for which the scale is small enough for viscosity to become important (that is, viscous forces become of the order of inertial ones). It is at these small scales where the dissipation of energy by viscous action finally takes place. The Reynolds number indicates at what scale this viscous dissipation occurs.
In physiology
Poiseuille's law on blood circulation in the body is dependent on laminar flow. In turbulent flow the flow rate is proportional to the square root of the pressure gradient, as opposed to its direct proportionality to pressure gradient in laminar flow.
Using the definition of the Reynolds number we can see that a large diameter with rapid flow, where the density of the blood is high, tends towards turbulence. Rapid changes in vessel diameter may lead to turbulent flow, for instance when a narrower vessel widens to a larger one. Furthermore, a bulge of atheroma may be the cause of turbulent flow, where audible turbulence may be detected with a stethoscope.
Complex systems
Reynolds number interpretation has been extended into the area of arbitrary complex systems. Such as financial flows, nonlinear networks, etc. In the latter case, an artificial viscosity is reduced to a nonlinear mechanism of energy distribution in complex network media. Reynolds number then represents a basic control parameter that expresses a balance between injected and dissipated energy flows for an open boundary system. It has been shown that Reynolds critical regime separates two types of phase space motion: accelerator (attractor) and decelerator. High Reynolds number leads to a chaotic regime transition only in frame of strange attractor model.
Relationship to other dimensionless parameters
There are many dimensionless numbers in fluid mechanics. The Reynolds number measures the ratio of advection and diffusion effects on structures in the velocity field, and is therefore closely related to Péclet numbers, which measure the ratio of these effects on other fields carried by the flow, for example, temperature and magnetic fields. Replacement of the kinematic viscosity in by the thermal or magnetic diffusivity results in respectively the thermal Péclet number and the magnetic Reynolds number. These are therefore related to by-products with ratios of diffusivities, namely the Prandtl number and magnetic Prandtl number.
| Physical sciences | Fluid mechanics | null |
23870096 | https://en.wikipedia.org/wiki/Negative-index%20metamaterial | Negative-index metamaterial | Negative-index metamaterial or negative-index material (NIM) is a metamaterial whose refractive index for an electromagnetic wave has a negative value over some frequency range.
NIMs are constructed of periodic basic parts called unit cells, which are usually significantly smaller than the wavelength of the externally applied electromagnetic radiation. The unit cells of the first experimentally investigated NIMs were constructed from circuit board material, or in other words, wires and dielectrics. In general, these artificially constructed cells are stacked or planar and configured in a particular repeated pattern to compose the individual NIM. For instance, the unit cells of the first NIMs were stacked horizontally and vertically, resulting in a pattern that was repeated and intended (see below images).
Specifications for the response of each unit cell are predetermined prior to construction and are based on the intended response of the entire, newly constructed, material. In other words, each cell is individually tuned to respond in a certain way, based on the desired output of the NIM. The aggregate response is mainly determined by each unit cell's geometry and substantially differs from the response of its constituent materials. In other words, the way the NIM responds is that of a new material, unlike the wires or metals and dielectrics it is made from. Hence, the NIM has become an effective medium. Also, in effect, this metamaterial has become an “ordered macroscopic material, synthesized from the bottom up”, and has emergent properties beyond its components.
Metamaterials that exhibit a negative value for the refractive index are often referred to by any of several terminologies: left-handed media or left-handed material (LHM), backward-wave media (BW media), media with negative refractive index, double negative (DNG) metamaterials, and other similar names.
Properties and characteristics
Electrodynamics of media with negative indices of refraction were first studied by Russian theoretical physicist Victor Veselago from Moscow Institute of Physics and Technology in 1967. The proposed left-handed or negative-index materials were theorized to exhibit optical properties opposite to those of glass, air, and other transparent media. Such materials were predicted to exhibit counterintuitive properties like bending or refracting light in unusual and unexpected ways. However, the first practical metamaterial was not constructed until 33 years later and it does support Veselago's concepts.
Currently, negative-index metamaterials are being developed to manipulate electromagnetic radiation in new ways. For example, optical and electromagnetic properties of natural materials are often altered through chemistry. With metamaterials, optical and electromagnetic properties can be engineered by changing the geometry of its unit cells. The unit cells are materials that are ordered in geometric arrangements with dimensions that are fractions of the wavelength of the radiated electromagnetic wave. Each artificial unit responds to the radiation from the source. The collective result is the material's response to the electromagnetic wave that is broader than normal.
Subsequently, transmission is altered by adjusting the shape, size, and configurations of the unit cells. This results in control over material parameters known as permittivity and magnetic permeability. These two parameters (or quantities) determine the propagation of electromagnetic waves in matter. Therefore, controlling the values of permittivity and permeability means that the refractive index can be negative or zero as well as conventionally positive. It all depends on the intended application or desired result. So, optical properties can be expanded beyond the capabilities of lenses, mirrors, and other conventional materials. Additionally, one of the effects most studied is the negative index of refraction.
Reverse propagation
When a negative index of refraction occurs, propagation of the electromagnetic wave is reversed. Resolution below the diffraction limit becomes possible. This is known as subwavelength imaging. Transmitting a beam of light via an electromagnetically flat surface is another capability. In contrast, conventional materials are usually curved, and cannot achieve resolution below the diffraction limit. Also, reversing the electromagnetic waves in a material, in conjunction with other ordinary materials (including air) could result in minimizing losses that would normally occur.
The reverse of the electromagnetic wave, characterized by an antiparallel phase velocity is also an indicator of negative index of refraction.
Furthermore, negative-index materials are customized composites. In other words, materials are combined with a desired result in mind. Combinations of materials can be designed to achieve optical properties not seen in nature. The properties of the composite material stem from its lattice structure constructed from components smaller than the impinging electromagnetic wavelength separated by distances that are also smaller than the impinging electromagnetic wavelength. Likewise, by fabricating such metamaterials researchers are trying to overcome fundamental limits tied to the wavelength of light. The unusual and counterintuitive properties currently have practical and commercial use manipulating electromagnetic microwaves in wireless and communication systems. Lastly, research continues in the other domains of the electromagnetic spectrum, including visible light.
Materials
The first actual metamaterials worked in the microwave regime, or centimeter wavelengths, of the electromagnetic spectrum (about 4.3 GHz). It was constructed of split-ring resonators and conducting straight wires (as unit cells). The unit cells were sized from 7 to 10 millimeters. The unit cells were arranged in a two-dimensional (periodic) repeating pattern which produces a crystal-like geometry. Both the unit cells and the lattice spacing were smaller than the radiated electromagnetic wave. This produced the first left-handed material when both the permittivity and permeability of the material were negative. This system relies on the resonant behavior of the unit cells. Below a group of researchers develop an idea for a left-handed metamaterial that does not rely on such resonant behavior.
Research in the microwave range continues with split-ring resonators and conducting wires. Research also continues in the shorter wavelengths with this configuration of materials and the unit cell sizes are scaled down. However, at around 200 terahertz issues arise which make using the split ring resonator problematic. "Alternative materials become more suitable for the terahertz and optical regimes." At these wavelengths selection of materials and size limitations become important. For example, in 2007 a 100 nanometer mesh wire design made of silver and woven in a repeating pattern transmitted beams at the 780 nanometer wavelength, the far end of the visible spectrum. The researchers believe this produced a negative refraction of 0.6. Nevertheless, this operates at only a single wavelength like its predecessor metamaterials in the microwave regime. Hence, the challenges are to fabricate metamaterials so that they "refract light at ever-smaller wavelengths" and to develop broad band capabilities.
Artificial transmission-line-media
In the metamaterial literature, medium or media refers to transmission medium or optical medium. In 2002, a group of researchers came up with the idea that in contrast to materials that depended on resonant behavior, non-resonant phenomena could surpass narrow bandwidth constraints of the wire/split-ring resonator configuration. This idea translated into a type of medium with broader bandwidth abilities, negative refraction, backward waves, and focusing beyond the diffraction limit.
They dispensed with split-ring-resonators and instead used a network of L–C loaded transmission lines. In metamaterial literature this became known as artificial transmission-line media. At that time it had the added advantage of being more compact than a unit made of wires and split ring resonators. The network was both scalable (from the megahertz to the tens of gigahertz range) and tunable. It also includes a method for focusing the wavelengths of interest.
By 2007 the negative refractive index transmission line was employed as a subwavelength focusing free-space flat lens. That this is a free-space lens is a significant advance. Part of prior research efforts targeted creating a lens that did not need to be embedded in a transmission line.
The optical domain
Metamaterial components shrink as research explores shorter wavelengths (higher frequencies) of the electromagnetic spectrum in the infrared and visible spectrums. For example, theory and experiment have investigated smaller horseshoe shaped split ring resonators designed with lithographic techniques, as well as paired metal nanorods or nanostrips, and nanoparticles as circuits designed with lumped element models
Applications
The science of negative-index materials is being matched with conventional devices that broadcast, transmit, shape, or receive electromagnetic signals that travel over cables, wires, or air. The materials, devices and systems that are involved with this work could have their properties altered or heightened. Hence, this is already happening with metamaterial antennas and related devices which are commercially available. Moreover, in the wireless domain these metamaterial apparatuses continue to be researched. Other applications are also being researched. These are electromagnetic absorbers such as radar-microwave absorbers, electrically small resonators, waveguides that can go beyond the diffraction limit, phase compensators, advancements in focusing devices (e.g. microwave lens), and improved electrically small antennas.
In the optical frequency regime developing the superlens may allow for imaging below the diffraction limit. Other potential applications for negative-index metamaterials are optical nanolithography, nanotechnology circuitry, as well as a near field superlens (Pendry, 2000) that could be useful for biomedical imaging and subwavelength photolithography.
Manipulating permittivity and permeability
To describe any electromagnetic properties of a given achiral material such as an optical lens, there are two significant parameters. These are permittivity, , and permeability, , which allow accurate prediction of light waves traveling within materials, and electromagnetic phenomena that occur at the interface between two materials.
For example, refraction is an electromagnetic phenomenon which occurs at the interface between two materials. Snell's law states that the relationship between the angle of incidence of a beam of electromagnetic radiation (light) and the resulting angle of refraction rests on the refractive indices, , of the two media (materials). The refractive index of an achiral medium is given by . Hence, it can be seen that the refractive index is dependent on these two parameters. Therefore, if designed or arbitrarily modified values can be inputs for and , then the behavior of propagating electromagnetic waves inside the material can be manipulated at will. This ability then allows for intentional determination of the refractive index.
For example, in 1967, Victor Veselago analytically determined that light will refract in the reverse direction (negatively) at the interface between a material with negative refractive index and a material exhibiting conventional positive refractive index. This extraordinary material was realized on paper with simultaneous negative values for and , and could therefore be termed a double negative material. However, in Veselago's day a material which exhibits double negative parameters simultaneously seemed impossible because no natural materials exist which can produce this effect. Therefore, his work was ignored for three decades. It was nominated for the Nobel Prize later.
In general the physical properties of natural materials cause limitations. Most dielectrics only have positive permittivities, > 0. Metals will exhibit negative permittivity, < 0 at optical frequencies, and plasmas exhibit negative permittivity values in certain frequency bands. Pendry et al. demonstrated that the plasma frequency can be made to occur in the lower microwave frequencies for metals with a material made of metal rods that replaces the bulk metal. However, in each of these cases permeability remains always positive. At microwave frequencies it is possible for negative μ to occur in some ferromagnetic materials. But the inherent drawback is they are difficult to find above terahertz frequencies. In any case, a natural material that can achieve negative values for permittivity and permeability simultaneously has not been found or discovered. Hence, all of this has led to constructing artificial composite materials known as metamaterials in order to achieve the desired results.
Negative index of refraction due to chirality
In case of chiral materials, the refractive index depends not only on permittivity and permeability , but also on the chirality parameter , resulting in distinct values for left and right circularly polarized waves, given by
A negative index will occur for waves of one circular polarization if > . In this case, it is not necessary that either or both and be negative to achieve a negative index of refraction. A negative refractive index due to chirality was predicted by Pendry and Tretyakov et al., and first observed simultaneously and independently by Plum et al. and Zhang et al. in 2009.
Physical properties never before produced in nature
Theoretical articles were published in 1996 and 1999 which showed that synthetic materials could be constructed to purposely exhibit a negative permittivity and permeability.
These papers, along with Veselago's 1967 theoretical analysis of the properties of negative-index materials, provided the background to fabricate a metamaterial with negative effective permittivity and permeability. See below.
A metamaterial developed to exhibit negative-index behavior is typically formed from individual components. Each component responds differently and independently to a radiated electromagnetic wave as it travels through the material. Since these components are smaller than the radiated wavelength it is understood that a macroscopic view includes an effective value for both permittivity and permeability.
Composite material
In the year 2000, David R. Smith's team of UCSD researchers produced a new class of composite materials by depositing a structure onto a circuit-board substrate consisting of a series of thin copper split-rings and ordinary wire segments strung parallel to the rings. This material exhibited unusual physical properties that had never been observed in nature. These materials obey the laws of physics, but behave differently from normal materials. In essence these negative-index metamaterials were noted for having the ability to reverse many of the physical properties that govern the behavior of ordinary optical materials. One of those unusual properties is the ability to reverse, for the first time, Snell's law of refraction. Until the demonstration of negative refractive index for microwaves by the UCSD team, the material had been unavailable. Advances during the 1990s in fabrication and computation abilities allowed these first metamaterials to be constructed. Thus, the "new" metamaterial was tested for the effects described by Victor Veselago 30 years earlier. Studies of this experiment, which followed shortly thereafter, announced that other effects had occurred.
With antiferromagnets and certain types of insulating ferromagnets, effective negative magnetic permeability is achievable when polariton resonance exists. To achieve a negative index of refraction, however, permittivity with negative values must occur within the same frequency range. The artificially fabricated split-ring resonator is a design that accomplishes this, along with the promise of dampening high losses. With this first introduction of the metamaterial, it appears that the losses incurred were smaller than antiferromagnetic, or ferromagnetic materials.
When first demonstrated in 2000, the composite material (NIM) was limited to transmitting microwave radiation at frequencies of 4 to 7 gigahertz (4.28–7.49 cm wavelengths). This range is between the frequency of household microwave ovens (~2.45 GHz, 12.23 cm) and military radars (~10 GHz, 3 cm). At demonstrated frequencies, pulses of electromagnetic radiation moving through the material in one direction are composed of constituent waves moving in the opposite direction.
The metamaterial was constructed as a periodic array of copper split ring and wire conducting elements deposited onto a circuit-board substrate. The design was such that the cells, and the lattice spacing between the cells, were much smaller than the radiated electromagnetic wavelength. Hence, it behaves as an effective medium. The material has become notable because its range of (effective) permittivity εeff and permeability μeff values have exceeded those found in any ordinary material. Furthermore, the characteristic of negative (effective) permeability evinced by this medium is particularly notable, because it has not been found in ordinary materials. In addition, the negative values for the magnetic component is directly related to its left-handed nomenclature, and properties (discussed in a section below). The split-ring resonator (SRR), based on the prior 1999 theoretical article, is the tool employed to achieve negative permeability. This first composite metamaterial is then composed of split-ring resonators and electrical conducting posts.
Initially, these materials were only demonstrated at wavelengths longer than those in the visible spectrum. In addition, early NIMs were fabricated from opaque materials and usually made of non-magnetic constituents. As an illustration, however, if these materials are constructed at visible frequencies, and a flashlight is shone onto the resulting NIM slab, the material should focus the light at a point on the other side. This is not possible with a sheet of ordinary opaque material. In 2007, the NIST in collaboration with the Atwater Lab at Caltech created the first NIM active at optical frequencies. More recently (), layered "fishnet" NIM materials made of silicon and silver wires have been integrated into optical fibers to create active optical elements.
Simultaneous negative permittivity and permeability
Negative permittivity εeff < 0 had already been discovered and realized in metals for frequencies all the way up to the plasma frequency, before the first metamaterial. There are two requirements to achieve a negative value for refraction. First, is to fabricate a material which can produce negative permeability μeff < 0. Second, negative values for both permittivity and permeability must occur simultaneously over a common range of frequencies.
Therefore, for the first metamaterial, the nuts and bolts are one split-ring resonator electromagnetically combined with one (electric) conducting post. These are designed to resonate at designated frequencies to achieve the desired values. Looking at the make-up of the split ring, the associated magnetic field pattern from the SRR is dipolar. This dipolar behavior is notable because this means it mimics nature's atom, but on a much larger scale, such as in this case at 2.5 millimeters. Atoms exist on the scale of picometers.
The splits in the rings create a dynamic where the SRR unit cell can be made resonant at radiated wavelengths much larger than the diameter of the rings. If the rings were closed, a half wavelength boundary would be electromagnetically imposed as a requirement for resonance.
The split in the second ring is oriented opposite to the split in the first ring. It is there to generate a large capacitance, which occurs in the small gap. This capacitance substantially decreases the resonant frequency while concentrating the electric field. The individual SRR depicted on the right had a resonant frequency of 4.845 GHz, and the resonance curve, inset in the graph, is also shown. The radiative losses from absorption and reflection are noted to be small, because the unit dimensions are much smaller than the free space, radiated wavelength.
When these units or cells are combined into a periodic arrangement, the magnetic coupling between the resonators is strengthened, and a strong magnetic coupling occurs. Properties unique in comparison to ordinary or conventional materials begin to emerge. For one thing, this periodic strong coupling creates a material, which now has an effective magnetic permeability μeff in response to the radiated-incident magnetic field.
Composite material passband
Graphing the general dispersion curve, a region of propagation occurs from zero up to a lower band edge, followed by a gap, and then an upper passband. The presence of a 400 MHz gap between 4.2 GHz and 4.6 GHz implies a band of frequencies where μeff < 0 occurs.
(Please see the image in the previous section)
Furthermore, when wires are added symmetrically between the split rings, a passband occurs within the previously forbidden band of the split ring dispersion curves. That this passband occurs within a previously forbidden region indicates that the negative εeff for this region has combined with the negative μeff to allow propagation, which fits with theoretical predictions. Mathematically, the dispersion relation leads to a band with negative group velocity everywhere, and a bandwidth that is independent of the plasma frequency, within the stated conditions.
Mathematical modeling and experiment have both shown that periodically arrayed conducting elements (non-magnetic by nature) respond predominantly to the magnetic component of incident electromagnetic fields. The result is an effective medium and negative μeff over a band of frequencies. The permeability was verified to be the region of the forbidden band, where the gap in propagation occurred – from a finite section of material. This was combined with a negative permittivity material, εeff < 0, to form a “left-handed” medium, which formed a propagation band with negative group velocity where previously there was only attenuation. This validated predictions. In addition, a later work determined that this first metamaterial had a range of frequencies over which the refractive index was predicted to be negative for one direction of propagation (see ref #). Other predicted electrodynamic effects were to be investigated in other research.
Describing a left-handed material
From the conclusions in the above section a left-handed material (LHM) can be defined. It is a material which exhibits simultaneous negative values for permittivity, ε, and permeability, μ, in an overlapping frequency region. Since the values are derived from the effects of the composite medium system as a whole, these are defined as effective permittivity, εeff, and effective permeability, μeff. Real values are then derived to denote the value of negative index of refraction, and wave vectors. This means that in practice losses will occur for a given medium used to transmit electromagnetic radiation such as microwave, or infrared frequencies, or visible light – for example. In this instance, real values describe either the amplitude or the intensity of a transmitted wave relative to an incident wave, while ignoring the negligible loss values.
Isotropic negative index in two dimensions
In the above sections first fabricated metamaterial was constructed with resonating elements, which exhibited one direction of incidence and polarization. In other words, this structure exhibited left-handed propagation in one dimension. This was discussed in relation to Veselago's seminal work 33 years earlier (1967). He predicted that intrinsic to a material, which manifests negative values of effective permittivity and permeability, are several types of reversed physics phenomena. Hence, there was then a critical need for a higher-dimensional LHMs to confirm Veselago's theory, as expected. The confirmation would include reversal of Snell's law (index of refraction), along with other reversed phenomena.
In the beginning of 2001 the existence of a higher-dimensional structure was reported. It was two-dimensional and demonstrated by both experiment and numerical confirmation. It was an LHM, a composite constructed of wire strips mounted behind the split-ring resonators (SRRs) in a periodic configuration. It was created for the express purpose of being suitable for further experiments to produce the effects predicted by Veselago.
Experimental verification of a negative index of refraction
A theoretical work published in 1967 by Soviet physicist Victor Veselago showed that a refractive index with negative values is possible and that this does not violate the laws of physics. As discussed previously (above), the first metamaterial had a range of frequencies over which the refractive index was predicted to be negative for one direction of propagation. It was reported in May 2000.
In 2001, a team of researchers constructed a prism composed of metamaterials (negative-index metamaterials) to experimentally test for negative refractive index. The experiment used a waveguide to help transmit the proper frequency and isolate the material. This test achieved its goal because it successfully verified a negative index of refraction.
The experimental demonstration of negative refractive index was followed by another demonstration, in 2003, of a reversal of Snell's law, or reversed refraction. However, in this experiment negative index of refraction material is in free space from 12.6 to 13.2 GHz. Although the radiated frequency range is about the same, a notable distinction is this experiment is conducted in free space rather than employing waveguides.
Furthering the authenticity of negative refraction, the power flow of a wave transmitted through a dispersive left-handed material was calculated and compared to a dispersive right-handed material. The transmission of an incident field, composed of many frequencies, from an isotropic nondispersive material into an isotropic dispersive media is employed. The direction of power flow for both nondispersive and dispersive media is determined by the time-averaged
Poynting vector. Negative refraction was shown to be possible for multiple frequency signals by explicit calculation of the Poynting vector in the LHM.
Fundamental electromagnetic properties of the NIM
In a slab of conventional material with an ordinary refractive index – a right-handed material (RHM) – the wave front is transmitted away from the source. In a NIM the wavefront travels toward the source. However, the magnitude and direction of the flow of energy essentially remains the same in both the ordinary material and the NIM. Since the flow of energy remains the same in both materials (media), the impedance of the NIM matches the RHM. Hence, the sign of the intrinsic impedance is still positive in a NIM.
Light incident on a left-handed material, or NIM, will bend to the same side as the incident beam, and for Snell's law to hold, the refraction angle should be negative. In a passive metamaterial medium this determines a negative real and imaginary part of the refractive index.
Negative refractive index in left-handed materials
In 1968 Victor Veselago's paper showed that the opposite directions of EM plane waves and the flow of energy was derived from the individual Maxwell curl equations. In ordinary optical materials, the curl equation for the electric field show a "right hand rule" for the directions of the electric field E, the magnetic induction B, and wave propagation, which goes in the direction of wave vector k. However, the direction of energy flow formed by E × H is right-handed only when permeability is greater than zero. This means that when permeability is less than zero, e.g. negative, wave propagation is reversed (determined by k), and contrary to the direction of energy flow. Furthermore, the relations of vectors E, H, and k form a "left-handed" system – and it was Veselago who coined the term "left-handed" (LH) material, which is in wide use today (2011). He contended that an LH material has a negative refractive index and relied on the steady-state solutions of Maxwell's equations as a center for his argument.
After a 30-year void, when LH materials were finally demonstrated, it could be said that the designation of negative refractive index is unique to LH systems; even when compared to photonic crystals. Photonic crystals, like many other known systems, can exhibit unusual propagation behavior such as reversal of phase and group velocities. But, negative refraction does not occur in these systems, and not yet realistically in photonic crystals.
Negative refraction at optical frequencies
The negative refractive index in the optical range was first demonstrated in 2005 by Shalaev et al. (at the telecom wavelength λ = 1.5 μm) and by Brueck et al. (at λ = 2 μm) at nearly the same time.
In 2006, a Caltech team led by Lezec, Dionne, and Atwater achieved negative refraction in the visible spectral regime.
Reversed Cherenkov radiation
Besides reversed values for the index of refraction, Veselago predicted the occurrence of reversed Cherenkov radiation in a left-handed medium. Whereas ordinary Cherenkov radiation is emitted in a cone around the direction in which a charged particle is travelling through the medium, reversed Cherenkov radiation is emitted in a cone around the opposite direction. Reversed Cherenkov radiation was first experimentally demonstrated indirectly in 2009, using a phased electromagnetic dipole array to model a moving charged particle. Reversed Cherenkov radiation emitted by actual charged particles was first observed in 2017.
Other optics with NIMs
Theoretical work, along with numerical simulations, began in the early 2000s on the abilities of DNG slabs for subwavelength focusing. The research began with Pendry's proposed "Perfect lens." Several research investigations that followed Pendry's concluded that the "Perfect lens" was possible in theory but impractical. One direction in subwavelength focusing proceeded with the use of negative-index metamaterials, but based on the enhancements for imaging with surface plasmons. In another direction researchers explored paraxial approximations of NIM slabs.
Implications of negative refractive materials
The existence of negative refractive materials can result in a change in electrodynamic calculations for the case of permeability μ = 1 . A change from a conventional refractive index to a negative value gives incorrect results for conventional calculations, because some properties and effects have been altered. When permeability μ has values other than 1 this affects Snell's law, the Doppler effect, the Cherenkov radiation, Fresnel's equations, and Fermat's principle.
The refractive index is basic to the science of optics. Shifting the refractive index to a negative value may be a cause to revisit or reconsider the interpretation of some norms, or basic laws.
US patent on left-handed composite media
The first US patent for a fabricated metamaterial, titled "Left handed composite media" by David R. Smith, Sheldon Schultz, Norman Kroll and Richard A. Shelby, was issued in 2004. The invention achieves simultaneous negative permittivity and permeability over a common band of frequencies. The material can integrate media which is already composite or continuous, but which will produce negative permittivity and permeability within the same spectrum of frequencies. Different types of continuous or composite may be deemed appropriate when combined for the desired effect. However, the inclusion of a periodic array of conducting elements is preferred. The array scatters electromagnetic radiation at wavelengths longer than the size of the element and lattice spacing. The array is then viewed as an effective medium.
| Physical sciences | Basics_3 | Physics |
23874520 | https://en.wikipedia.org/wiki/Velocity%20dispersion | Velocity dispersion | In astronomy, the velocity dispersion (σ) is the statistical dispersion of velocities about the mean velocity for a group of astronomical objects, such as an open cluster, globular cluster, galaxy, galaxy cluster, or supercluster. By measuring the radial velocities of the group's members through astronomical spectroscopy, the velocity dispersion of that group can be estimated and used to derive the group's mass from the virial theorem. Radial velocity is found by measuring the Doppler width of spectral lines of a collection of objects; the more radial velocities one measures, the more accurately one knows their dispersion. A central velocity dispersion refers to the σ of the interior regions of an extended object, such as a galaxy or cluster.
The relationship between velocity dispersion and matter (or the observed electromagnetic radiation emitted by this matter) takes several forms – specific correlations – in astronomy based on the object(s) being observed. Notably, the M–σ relation applies for material orbiting many black holes, the Faber–Jackson relation for elliptical galaxies, and the Tully–Fisher relation for spiral galaxies. For example, the σ found for objects about the Milky Way's supermassive black hole (SMBH) is about 100 km/s, which provides an approximation of the mass of this SMBH. The Andromeda Galaxy (Messier 31) hosts a SMBH about 10 times larger than our own, and has a .
Groups and clusters of galaxies have more disparate (contrasting in degree) velocity dispersions than smaller objects. For example, while our own poor group, the Local Group, has a , rich clusters of galaxies, such as the Coma Cluster, have a . The dwarf elliptical galaxies within Coma, as with all galaxies, have their own internal velocity dispersion for their stars, which is a , typically. Normal elliptical galaxies, by comparison, have an average .
For spiral galaxies, the increase in velocity dispersion in population I stars is a gradual process which likely results from the near-random incidence of momentum exchanges, specifically dynamical friction, between individual stars and large interstellar media (gas and dust clouds) with masses greater than . Face-on spiral galaxies have a central ; slightly more if viewed edge-on.
| Physical sciences | Celestial mechanics | Astronomy |
20865567 | https://en.wikipedia.org/wiki/Schiehallion%20experiment | Schiehallion experiment | The Schiehallion experiment was an 18th-century experiment to determine the mean density of the Earth. Funded by a grant from the Royal Society, it was conducted in the summer of 1774 around the Scottish mountain of Schiehallion, Perthshire. The experiment involved measuring the tiny deflection of the vertical due to the gravitational attraction of a nearby mountain. Schiehallion was considered the ideal location after a search for candidate mountains, thanks to its isolation and almost symmetrical shape.
The experiment had previously been considered, but rejected, by Isaac Newton as a practical demonstration of his theory of gravitation; however, a team of scientists, notably Nevil Maskelyne, the Astronomer Royal, was convinced that the effect would be detectable and undertook to conduct the experiment. The deflection angle depended on the relative densities and volumes of the Earth and the mountain: if the density and volume of Schiehallion could be ascertained, then so could the density of the Earth. Once this was known, it would in turn yield approximate values for those of the other planets, their moons, and the Sun, previously known only in terms of their relative ratios.
Background
A pendulum hangs straight downwards in a symmetrical gravitational field. However, if a sufficiently large mass such as a mountain is nearby, its gravitational attraction should pull the pendulum's plumb-bob slightly out of true (in the sense that it doesn't point to the centre of mass of the Earth). The change in plumb-line angle against a known object—such as a star—could be carefully measured on opposite sides of the mountain. If the mass of the mountain could be independently established from a determination of its volume and an estimate of the mean density of its rocks, then these values could be extrapolated to provide the mean density of the Earth, and by extension, its mass.
Isaac Newton had considered the effect in the Principia, but pessimistically thought that any real mountain would produce too small a deflection to measure. Gravitational effects, he wrote, were only discernible on the planetary scale. Newton's pessimism was unfounded: although his calculations had suggested a deviation of less than 2 minutes of arc (for an idealised mountain), this angle, though very slight, was within the theoretical capability of instruments of his day.
An experiment to test Newton's idea would provide supporting evidence for both his law of universal gravitation and estimates of the mass and density of the Earth. Since the masses of astronomical objects were known only in terms of relative ratios, the mass of the Earth would provide reasonable values to the other planets, their moons, and the Sun. The data were also capable of determining the value of the Newtonian constant of gravitation , though this was not a goal of the experimenters; references to a value for would not appear in the scientific literature until almost a hundred years later.
Finding the mountain
Chimborazo, 1738
A pair of French astronomers, Pierre Bouguer and Charles Marie de La Condamine, were the first to attempt the experiment, conducting their measurements on the volcano Chimborazo. At the time, this lay in the "Real Audiencia of Quito" of the Viceroyalty of Peru, and is now in the province of Chimborazo in the Republic of Ecuador. Their expedition had left France for South America in 1735 to try to measure the meridian arc length of one degree of latitude near the equator, but they took advantage of the opportunity to attempt the deflection experiment. In December 1738, under very difficult conditions of terrain and climate, they conducted a pair of measurements at altitudes of 4,680 and 4,340 m. Bouguer wrote in a 1749 paper that they had been able to detect a deflection of 8 seconds of arc, but he downplayed the significance of their results, suggesting that the experiment would be better carried out under easier conditions in France or England. He added that the experiment had at least proved that the Earth could not be a hollow shell, as some thinkers of the day, including Edmond Halley, had suggested.
Schiehallion, 1774
Between 1763 and 1767, during operations to survey the Mason–Dixon line between the states of Pennsylvania and Maryland, British astronomers found many more systematic and non-random errors than might have been expected, extending the work longer than planned. When this information reached the members of the Royal Society, Henry Cavendish realized that this effect might have been due to the gravitational pull of the nearby Allegheny Mountains, which had probably diverted the plumb lines of the theodolites and the liquids inside spirit levels.
Prompted by this news, a further attempt on the experiment was proposed to the Royal Society in 1772 by Nevil Maskelyne, Astronomer Royal. He suggested that the experiment would "do honour to the nation where it was made" and proposed Whernside in Yorkshire, or the Blencathra-Skiddaw massif in Cumberland, as suitable targets. The Royal Society formed the Committee of Attraction to consider the matter, appointing Maskelyne, Joseph Banks and Benjamin Franklin amongst its members. The Committee dispatched the astronomer and surveyor Charles Mason to find a suitable mountain.
After a lengthy search over the summer of 1773, Mason reported that the best candidate was Schiehallion (then spelled Schehallien), a peak lying between Loch Tay and Loch Rannoch in the central Scottish Highlands. The mountain stood apart from any nearby hills, which would reduce their gravitational influence, and its symmetrical east–west ridge would simplify the calculations. Its steep northern and southern slopes would allow the experiment to be sited close to its centre of mass, maximising the deflection effect.
Mason declined to conduct the work himself for the offered commission of one guinea per day. The task therefore fell to Maskelyne, for which he was granted a temporary leave of his duties as Astronomer Royal. He was aided in the task by mathematician and surveyor Charles Hutton, and Reuben Burrow who was a mathematician of the Royal Greenwich Observatory. A workforce of labourers was engaged to construct observatories for the astronomers and assist in the surveying. The science team was particularly well-equipped: its astronomical instruments included a brass quadrant from Cook's 1769 transit of Venus expedition, a zenith sector, and a regulator (precision pendulum clock) for timing the astronomical observations. They also acquired a theodolite and Gunter's chain for surveying the mountain, and a pair of barometers for measuring altitude. Generous funding for the experiment was available due to underspend on the transit of Venus expedition, which had been turned over to the Society by King George III.
Measurements
Astronomical
Observatories were constructed to the north and south of the mountain, plus a bothy to accommodate equipment and the scientists. The ruins of these structures remain on the mountainside. Most of the workforce was housed in rough canvas tents. Maskelyne's astronomical measurements were the first to be conducted. It was necessary for him to determine the zenith distances with respect to the plumb line for a set of stars at the precise time that each passed due south (astronomic latitude). Weather conditions were frequently unfavourable due to mist and rain. However, from the south observatory, he was able to take 76 measurements on 34 stars in one direction, and then 93 observations on 39 stars in the other. From the north side, he then conducted a set of 68 observations on 32 stars and a set of 100 on 37 stars. By conducting sets of measurements with the plane of the zenith sector first facing east and then west, he successfully avoided any systematic errors arising from collimating the sector.
To determine the deflection due to the mountain, it was necessary to account for the curvature of the Earth: an observer moving north or south will see the local zenith shift by the same angle as any change in geodetic latitude. After accounting for observational effects such as precession, aberration of light and nutation, Maskelyne showed that the difference between the locally determined zenith for observers north and south of Schiehallion was 54.6 arc seconds. Once the surveying team had provided a difference of 42.94″ latitude between the two stations, he was able to subtract this, and after rounding to the accuracy of his observations, announce that the sum of the north and south deflections was 11.6″.
Maskelyne published his initial results in the Philosophical Transactions of the Royal Society in 1775, using preliminary data on the mountain's shape and hence the position of its center of gravity. This led him to expect a deflection of 20.9″ if the mean densities of Schiehallion and the Earth were equal. Since the deflection was about half this, he was able to make a preliminary announcement that the mean density of the Earth was approximately double that of Schiehallion. A more accurate value would have to await completion of the surveying process.
Maskelyne took the opportunity to note that Schiehallion exhibited a gravitational attraction, and thus all mountains did; and that Newton's inverse square law of gravitation had been confirmed. An appreciative Royal Society presented Maskelyne with the 1775 Copley Medal. The biographer Chalmers later noted that "If any doubts yet remained with respect to the truth of the Newtonian system, they were now totally removed".
Surveying
The work of the surveying team was greatly hampered by the inclemency of the weather, and it took until 1776 to complete the task. To find the volume of the mountain, it was necessary to divide it into a set of vertical prisms and compute the volume of each. The triangulation task falling to Charles Hutton was considerable: the surveyors had obtained thousands of bearing angles to more than a thousand points around the mountain. Moreover, the vertices of his prisms did not always conveniently coincide with the surveyed heights. To make sense of all his data, he hit upon the idea of interpolating a series of lines at set intervals between his measured values, marking points of equal height. In doing so, not only could he easily determine the heights of his prisms, but from the swirl of the lines one could get an instant impression of the form of the terrain. Hutton thus used contour lines, which later became commonly used for depicting cartographic relief.
Hutton had to compute the individual attractions due to each of the many prisms that formed his grid, a process which was as laborious as the survey itself. The task occupied his time for a further two years before he could present his results, which he did in a hundred-page paper to the Royal Society in 1778. From his calculations, which took into account the effect of latitude on gravity, it followed that if the density of the Earth and Schiehallion had been the same, the attraction of the plumb-bob to the Earth would be 9,933 times the sum of its attractions to the mountain at the north and south stations. Yet the measured deflection of 11.6″ meant that Earth attraction was actually 17,804 times as great. From this it follows that the average density of Earth is approximately 1.8 times the density of the mountain. Hutton took a density of for Schiehallion, and announced that the density of the Earth was 1.8 times this, or , less than 20% away from the modern value of .
That the mean density of the Earth should so greatly exceed that of its surface rocks naturally meant that there must be more dense material lying deeper. Hutton correctly surmised that the core material was likely metallic, and might have a density of . He estimated this metallic portion to occupy some 65% of the diameter of the Earth. With a value for the mean density of the Earth, Hutton was able to set some values to Jérôme Lalande's planetary tables, which had previously only been able to express the densities of the major solar system objects in relative terms.
Repeat experiments
A more accurate measurement of the mean density of the Earth was made 24 years after Schiehallion, when in 1798 Henry Cavendish used an exquisitely sensitive torsion balance to measure the attraction between large masses of lead. Cavendish's figure of was only 1.2% from the currently accepted value of , and his result would not be significantly improved upon until 1895 by Charles Boys. The care with which Cavendish conducted the experiment and the accuracy of his result has led his name to since be associated with it.
The Scottish scientist John Playfair carried out a second survey of Schiehallion in 1811; on the basis of a rethink of its rock strata, he suggested a density of 4,560 to , though the then elderly Hutton vigorously defended the original value in an 1821 paper to the Society. Playfair's calculations had raised the density closer towards its modern value, but was still too low and significantly poorer than Cavendish's computation of some years earlier.
The Schiehallion experiment was repeated in 1856 by Henry James, director-general of the Ordnance Survey, who instead used the hill Arthur's Seat in central Edinburgh. With the resources of the Ordnance Survey at his disposal, James extended his topographical survey to a 21-kilometre radius, taking him as far as the borders of Midlothian. He obtained a density of about .
An experiment in 2005 undertook a variation of the 1774 work: instead of computing local differences in the zenith, the experiment made a very accurate comparison of the period of a pendulum at the top and bottom of Schiehallion. The period of a pendulum is a function of g, the local gravitational acceleration. The pendulum is expected to run more slowly at altitude, but the mass of the mountain will act to reduce this difference. This experiment has the advantage of being considerably easier to conduct than the 1774 one, but to achieve the desired accuracy, it is necessary to measure the period of the pendulum to within one part in one million. This experiment yielded a value of the mass of the Earth of , corresponding to a mean density of .
A modern re-examination of the geophysical data was able to take account of factors the 1774 team could not. With the benefit of a 120-km radius digital elevation model, greatly improved knowledge of the geology of Schiehallion, and the help of a computer, a 2007 report produced a mean Earth density of . When compared to the modern figure of , it stood as a testament to the accuracy of Maskelyne's astronomical observations.
Mathematical procedure
Consider the force diagram to the right, in which the deflection has been greatly exaggerated. The analysis has been simplified by considering the attraction on only one side of the mountain. A plumb-bob of mass is situated a distance from , the centre of mass of a mountain of mass and density . It is deflected through a small angle due to its attraction towards and its weight directed towards the Earth. The vector sum of and results in a tension in the pendulum string. The Earth has a mass , radius and a density .
The two gravitational forces on the plumb-bob are given by Newton's law of gravitation:
where is the Newtonian constant of gravitation. and can be eliminated by taking the ratio of to :
where and are the volumes of the mountain and the Earth. Under static equilibrium, the horizontal and vertical components of the string tension can be related to the gravitational forces and the deflection angle :
Substituting for :
Since , and are all known, has been measured and has been computed, then a value for the ratio can be obtained:
| Physical sciences | Earth science basics: General | Earth science |
20869694 | https://en.wikipedia.org/wiki/Atrial%20fibrillation | Atrial fibrillation | Atrial fibrillation (AF, AFib or A-fib) is an abnormal heart rhythm (arrhythmia) characterized by rapid and irregular beating of the atrial chambers of the heart. It often begins as short periods of abnormal beating, which become longer or continuous over time. It may also start as other forms of arrhythmia such as atrial flutter that then transform into AF.
Episodes can be asymptomatic. Symptomatic episodes may involve heart palpitations, fainting, lightheadedness, loss of consciousness, shortness of breath, or chest pain. Atrial fibrillation is associated with an increased risk of heart failure, dementia, and stroke. It is a type of supraventricular tachycardia.
Atrial fibrillation frequently results from bursts of tachycardia that originate in muscle bundles extending from the atrium to the pulmonary veins. Pulmonary vein isolation by transcatheter ablation can restore sinus rhythm. The ganglionated plexi (autonomic ganglia of the heart atrium and ventricles) can also be a source of atrial fibrillation, and is sometimes also ablated for that reason. Not only the pulmonary vein, but the left atrial appendage can be a source of atrial fibrillation and is also ablated for that reason. As atrial fibrillation becomes more persistent, the junction between the pulmonary veins and the left atrium becomes less of an initiator and the left atrium becomes an independent source of arrhythmias.
High blood pressure and valvular heart disease are the most common modifiable risk factors for AF. Other heart-related risk factors include heart failure, coronary artery disease, cardiomyopathy, and congenital heart disease. In low- and middle-income countries, valvular heart disease is often attributable to rheumatic fever. Lung-related risk factors include COPD, obesity, and sleep apnea. Cortisol and other stress biomarkers (including vasopressin, chromogranin A, and heat shock proteins), as well as emotional stress, may play a role in the pathogenesis of atrial fibrillation.
Other risk factors include excess alcohol intake, tobacco smoking, diabetes mellitus, and thyrotoxicosis. However, about half of cases are not associated with any of these aforementioned risks. Moreover, thyrotoxicosis seems to be an especially rare risk factor. Healthcare professionals might suspect AF after feeling the pulse and confirm the diagnosis by interpreting an electrocardiogram (ECG). A typical ECG in AF shows irregularly spaced QRS complexes without P waves.
Healthy lifestyle changes, such as weight loss in people with obesity, increased physical activity, and drinking less alcohol, can lower the risk for AF and reduce its burden if it occurs. AF is often treated with medications to slow the heart rate to a near-normal range (known as rate control) or to convert the rhythm to normal sinus rhythm (known as rhythm control). Electrical cardioversion can convert AF to normal heart rhythm and is often necessary for emergency use if the person is unstable. Ablation may prevent recurrence in some people. For those at low risk of stroke, AF does not necessarily require blood-thinning though some healthcare providers may prescribe an anti-clotting medication. Most people with AF are at higher risk of stroke. For those at more than low risk, experts generally recommend an anti-clotting medication. Anti-clotting medications include warfarin and direct oral anticoagulants. While these medications reduce stroke risk, they increase rates of major bleeding.
Atrial fibrillation is the most common serious abnormal heart rhythm and, as of 2020, affects more than 33 million people worldwide. As of 2014, it affected about 2 to 3% of the population of Europe and North America. This was an increase from 0.4 to 1% of the population around 2005. In the developing world, about 0.6% of males and 0.4% of females are affected. The percentage of people with AF increases with age with 0.1% under 50 years old, 4% between 60 and 70 years old, and 14% over 80 years old being affected. A-fib and atrial flutter resulted in 193,300 deaths in 2015, up from 29,000 in 1990. The first known report of an irregular pulse was by Jean-Baptiste de Sénac in 1749. Thomas Lewis was the first doctor to document this by ECG in 1909.
Signs and symptoms
AF is usually accompanied by symptoms related to a rapid heart rate. Rapid and irregular heart rates may be perceived as the sensation of the heart beating too fast, irregularly, or skipping beats (palpitations) or exercise intolerance and occasionally may produce anginal chest pain (if the high heart rate causes the heart's demand for oxygen to increase beyond the supply of available oxygen). Other possible symptoms include congestive heart failure symptoms such as fatigue, shortness of breath, or swelling. Loss of consciousness can also occur on atrial fibrillations due to lack of oxygen and blood to the brain. The abnormal heart rhythm (arrhythmia) is sometimes only identified with the onset of a stroke or a transient ischemic attack (TIA). It is not uncommon for a person to first become aware of AF from a routine physical examination or electrocardiogram, as it often does not cause symptoms.
Since most cases of AF are secondary to other medical problems, the presence of chest pain or angina, signs and symptoms of hyperthyroidism (an overactive thyroid gland) such as weight loss and diarrhea, and symptoms suggestive of lung disease can indicate an underlying cause. A history of stroke or TIA, as well as high blood pressure, diabetes, heart failure, or rheumatic fever, may indicate whether someone with AF is at a higher risk of complications.
Rapid heart rate
Presentation is similar to other forms of rapid heart rate and may be asymptomatic. Palpitations and chest discomfort are common complaints. The rapid uncoordinated heart rate may result in reduced output of blood pumped by the heart (cardiac output), resulting in inadequate blood flow, and therefore oxygen delivery to the rest of the body. Common symptoms of uncontrolled atrial fibrillation may include shortness of breath, shortness of breath when lying flat, dizziness, and sudden onset of shortness of breath during the night. This may progress to swelling of the lower extremities, a manifestation of congestive heart failure. Due to inadequate cardiac output, individuals with AF may also complain of lightheadedness.
AF can cause respiratory distress due to congestion in the lungs. By definition, the heart rate will be greater than 100 beats per minute. Blood pressure may be variable, and often difficult to measure as the beat-by-beat variability causes problems for most digital (oscillometric) non-invasive blood pressure monitors. For this reason, when determining the heart rate in AF, direct cardiac auscultation is recommended. Low blood pressure is most concerning, and a sign that immediate treatment is required. Many of the symptoms associated with uncontrolled atrial fibrillation are a manifestation of congestive heart failure due to the reduced cardiac output. The affected person's respiratory rate often increases in the presence of respiratory distress. Pulse oximetry may confirm the presence of too little oxygen reaching the body's tissues, related to any precipitating factors such as pneumonia. Examination of the jugular veins may reveal elevated pressure (jugular venous distention). Examination of the lungs may reveal crackles, which are suggestive of pulmonary edema. Examination of the heart will reveal a rapid irregular rhythm.
Causes
AF is linked to several forms of cardiovascular disease but may occur in otherwise normal hearts. Cardiovascular factors known to be associated with the development of AF include high blood pressure, coronary artery disease, mitral valve stenosis (e.g., due to rheumatic heart disease or mitral valve prolapse), mitral regurgitation, left atrial enlargement, hypertrophic cardiomyopathy, pericarditis, congenital heart disease, and previous heart surgery. People with congenital heart disease tend to develop atrial fibrillation at a younger age, that is more likely to be of right atrial origin (atypical) than of left origin, and have a greater risk of progressing to permanent atrial fibrillation.
Additionally, lung diseases (such as pneumonia, lung cancer, pulmonary embolism, and sarcoidosis) may play a role in certain people. Sepsis also increases the risk of developing new-onset atrial fibrillation. Disorders of breathing during sleep, such as obstructive sleep apnea (OSA), are also associated with AF. OSA, specifically, was found to be a very strong predictor of atrial fibrillation. Patients with OSA were shown to have an increased incidence of atrial fibrillation and a study done by Gami et al. demonstrated that increased nocturnal oxygen desaturation from OSA severity was correlated with higher incidences of atrial fibrillation. Obesity is a risk factor for AF. Hyperthyroidism and subclinical hyperthyroidism are associated with AF development.
Caffeine consumption does not appear to be associated with AF; excessive alcohol consumption ("binge drinking" or "holiday heart syndrome") is linked to AF. Low-to-moderate alcohol consumption also appears to be associated with an increased risk of developing atrial fibrillation, although the increase in risk associated with drinking less than two drinks daily appears to be small. Tobacco smoking and secondhand tobacco smoke exposure are associated with an increased risk of developing atrial fibrillation. Long-term endurance exercise that far exceeds the recommended amount of exercise (e.g., long-distance cycling or marathon running) appears to be associated with a modest increase in the risk of atrial fibrillation in middle-aged and elderly people.
Major stress biomarkers (including cortisol and heat shock proteins) indicate that stress plays a significant role in causing atrial fibrillation. There is some evidence that night shift working may be linked to a diagnosis of AF.
Atrial fibrillation is associated with elevated levels of inflammatory markers and clotting factors. Mendelian randomization indicates a causal relationship of inflammation leading to atrial fibrillation.
Genetics
Family history in a first degree relative is associated with a 40% increase in risk of AF. This finding led to the mapping of different loci such as 10q22-24, 6q14-16 and 11p15-5.3 and discover mutations associated with the loci. Mutations have been found in the genes of K+ channels and Na+ channels which affect the processes of polarization-depolarization of the myocardium, cellular hyper-excitability, shortening of effective refractory period favoring re-entries.
Using genome-wide association study (GWAS), which screen the entire genome for single nucleotide polymorphism (SNP), three susceptibility loci have been found for AF (4q25, 1q21 and 16q22). In these loci there are SNPs associated with a 30% increase in risk of recurrent atrial tachycardia after ablation. There are also SNPs associated with loss of function of the Pitx2c gene (involved in cellular development of pulmonary valves), responsible for re-entries. There are also SNPs close to ZFHX3 genes involved in the regulation of Ca2+. A 2018 meta-analysis of GWAS studies identified 97 locis associated with AF, of which 70 were newly identified associations: they are associated with genes that encode transcription factors, such as TBX3 and TBX5, NKX2-5 or PITX2, involved in the regulation of cardiac conduction, modulation of ion channels and in cardiac development.
Sedentary lifestyle
A sedentary lifestyle increases the risk factors associated with AF, such as obesity, hypertension, or diabetes mellitus. This favors remodeling processes of the atrium due to inflammation or alterations in the depolarization of cardiomyocytes by elevation of sympathetic nervous system activity. A sedentary lifestyle is associated with an increased risk of AF compared to physical activity. In both men and women, the practice of moderate exercise reduces the risk of AF progressively; intense sports may increase the risk of developing AF, as seen in athletes. It is due to a remodeling of cardiac tissue, and an increase in vagal tone, which shortens the effective refractory period (ERP) favoring re-entries from the pulmonary veins.
Tobacco
The rate of AF in smokers is 1.4 times higher than in non-smokers. Snus consumption, which delivers nicotine at a dose equivalent to that of cigarettes, is not correlated with AF.
Alcohol
Acute alcohol consumption can directly trigger an episode of atrial fibrillation. Regular alcohol consumption also increases the risk of atrial fibrillation in several ways. The long-term use of alcohol alters the physical structure and electrical properties of the atria. Alcohol consumption does this by repeatedly stimulating the sympathetic nervous system, increasing inflammation in the atria, raising blood pressure, lowering the levels of potassium and magnesium in the blood, worsening obstructive sleep apnea, and by promoting harmful structural changes (remodeling) in the atria and ventricles of the heart. This remodeling leads to abnormally increased pressure in the left atrium, inappropriately dilates it, and increases scarring (fibrosis) in the left atrium. The aforementioned structural changes increase the risk of developing atrial fibrillation when paired with the harmful changes in how the left atrium conducts electricity.
High blood pressure (hypertension)
In patients with hypertension prevalence rates reportedly range from 49% to 90%. According to the CHARGE Consortium, both systolic and diastolic blood pressure are predictors of the risk of AF. Systolic blood pressure values close to normal limit the increase in the risk associated with AF. Diastolic dysfunction is also associated with AF, which increases left atrial pressure, left atrial volume, size, and left ventricular hypertrophy, characteristic of chronic hypertension. All atrial remodeling is related to heterogeneous conduction and the formation of re-entrant electric conduction from the pulmonary veins.
Other diseases
There is a relationship between risk factors such as obesity and hypertension, with the appearance of diseases such as diabetes mellitus and sleep apnea-hypopnea syndrome, specifically, obstructive sleep apnea (OSA). These diseases are associated with an increased risk of AF due to their remodeling effects on the left atrium.
Medications
Several medications are associated with an increased risk of developing atrial fibrillation. Few studies have examined this phenomenon, and the exact incidence of medication-induced atrial fibrillation is unknown. Medications that are commonly associated with an increased risk of developing atrial fibrillation include dobutamine and the chemotherapy agent cisplatin. Agents associated with a moderately increased risk include nonsteroidal anti-inflammatory drugs (e.g., ibuprofen), bisphosphonates, and other chemotherapeutic agents such as melphalan, interleukin 2, and anthracyclines. Other medications that rarely increase the risk of developing atrial fibrillation include adenosine, aminophylline, corticosteroids, ivabradine, ondansetron, and antipsychotics. This form of atrial fibrillation occurs in people of all ages but is most common in the elderly, in those with other atrial fibrillation risk factors, and after heart surgery.
Pathophysiology
The normal electrical conduction system of the heart allows electrical impulses generated by the heart's own pacemaker (the sinoatrial node) to spread to and stimulate the muscular layer of the heart (myocardium) in both the atria and the ventricles. When the myocardium is stimulated it contracts, and if this occurs in an orderly manner allows blood to be pumped to the body. In AF, the normal regular electrical impulses generated by the sinoatrial node are overwhelmed by disorganized electrical waves, usually originating from the roots of the pulmonary veins. These disorganized waves conduct intermittently through the atrioventricular node, leading to irregular activation of the ventricles that generate the heartbeat.
Pathology
The primary pathologic change seen in atrial fibrillation is the progressive fibrosis of the atria. This fibrosis is due primarily to atrial dilation; however, genetic causes and inflammation may be factors in some individuals. Dilation of the atria can be due to almost any structural abnormality of the heart that can cause a rise in the pressure within the heart. This includes valvular heart disease (such as mitral stenosis, mitral regurgitation, and tricuspid regurgitation), hypertension, and congestive heart failure. Any inflammatory state that affects the heart can cause fibrosis of the atria.
Once dilation of the atria has occurred, this begins a chain of events that leads to the activation of the renin–angiotensin–aldosterone system (RAAS) and subsequent increase in the matrix metalloproteinases and disintegrin, which leads to atrial remodeling and fibrosis, with loss of atrial muscle mass. This process occurs gradually, and experimental studies have revealed patchy atrial fibrosis may precede the occurrence of atrial fibrillation and may progress with prolonged durations of atrial fibrillation.
Fibrosis is not limited to the muscle mass of the atria and may occur in the sinus node (SA node) and atrioventricular node (AV node), correlating with sick sinus syndrome. Prolonged episodes of atrial fibrillation have been shown to correlate with prolongation of the sinus node recovery time; this suggests that dysfunction of the SA node is progressive with prolonged episodes of atrial fibrillation.
Along with fibrosis, alterations in the atria that predispose to atrial fibrillation affect their electrical properties, as well as their responsiveness to the autonomic nervous system. The atrial remodeling that includes the pathologic changes described above has been referred to as atrial myopathy.
Electrophysiology
There are multiple theories about the cause of atrial fibrillation. An important theory is that the regular impulses produced by the sinus node for a normal heartbeat are overwhelmed by rapid electrical discharges produced in the atria and adjacent parts of the pulmonary veins. Non-pulmonary vein sources of triggers for atrial fibrillation have been identified in 10% to 33% of patients. These triggers include the coronary sinus, the posterior wall of the left atrium, and the left atrial appendage.
Sources of these disturbances are either automatic foci, often localized at one of the pulmonary veins, or a small number of localized sources in the form of either a re-entrant leading circle or electrical spiral waves (rotors); these localized sources may be in the left atrium near the pulmonary veins or in a variety of other locations through both the left or right atrium. Three fundamental components favor the establishment of a leading circle or a rotor: slow conduction velocity of the cardiac action potential, a short refractory period, and a small wavelength. Meanwhile, the wavelength is the product of velocity and refractory period. If the action potential has fast conduction, with a long refractory period and/or conduction pathway shorter than the wavelength, an AF focus would not be established. In multiple wavelet theory, a wavefront will break into smaller daughter wavelets when encountering an obstacle, through a process called vortex shedding. But, under the proper conditions, such wavelets can reform and spin around a center, forming an AF focus.
In a heart with AF, the increased calcium release from the sarcoplasmic reticulum and increased calcium sensitivity can lead to an accumulation of intracellular calcium and causes downregulation of L-type calcium channels. This reduces the duration of action potential and the refractory period, thus favoring the conduction of re-entrant waves. Increased expression of inward-rectifier potassium ion channels can cause a reduced atrial refractory period and wavelength. The abnormal distribution of gap junction proteins such as GJA1 (also known as connexin 43), and GJA5 (connexin 40) causes non-uniformity of electrical conduction, thus causing the arrhythmia.
AF can be distinguished from atrial flutter (AFL), which appears as an organized electrical circuit usually in the right atrium. AFL produces characteristic saw-toothed F-waves of constant amplitude and frequency on an ECG, whereas AF does not. In AFL, the discharges circulate rapidly at a rate of 300 beats per minute (bpm) around the atrium. In AF, there is no such regularity, except at the sources where the local activation rate can exceed 500 bpm. Although AF and atrial flutter are distinct arrhythmias, atrial flutter may degenerate into AF, and an individual may experience both arrhythmias at different times.
Although the electrical impulses of AF occur at a high rate, most of them do not result in a heartbeat. A heartbeat results when an electrical impulse from the atria passes through the atrioventricular (AV) node to the ventricles and causes them to contract. During AF, if all of the impulses from the atria passed through the AV node, there would be severe ventricular tachycardia, resulting in a severe reduction of cardiac output. This dangerous situation is prevented by the AV node since its limited conduction velocity reduces the rate at which impulses reach the ventricles during AF.
Diagnosis
Atrial fibrillation is diagnosed on an electrocardiogram (ECG/EKG).
The evaluation of atrial fibrillation involves a determination of the cause of the arrhythmia, and classification of the arrhythmia. Diagnostic investigation of AF typically includes a complete history and physical examination, ECG, transthoracic echocardiogram, complete blood count, serum thyroid stimulating hormone level and may include a functionality of some smartwatches.
Screening
Numerous guidelines recommend opportunistic screening for atrial fibrillation in those 65 years and older. These organizations include the: European Society of Cardiology,
National Heart Foundation of Australia and the Cardiac Society of Australia and New Zealand,
European Heart Rhythm Society,
AF-SCREEN International Collaboration,
Royal College of Physicians of Edinburgh
European Primary Care Cardiovascular Society, and Irish Health Information and Quality Authority.
Single timepoint screening detects undiagnosed AF, which is often asymptomatic, in approximately 1.4% of people in this age group. A Scottish inquiry into atrial fibrillation estimated that as many as one-third of people with AF are undiagnosed. Despite this, in 2018, the United States Preventive Services Task Force found insufficient evidence to determine the usefulness of routine screening.
Given the importance of having a pathway to treatment, general practice is potentially an ideal setting to conduct AF screening. General practice was identified as a 'preferred' setting for AF screening by the AF-SCREEN international collaboration report due to the availability of nursing support and the natural pathway to treatment.
Bloodwork
While many cases of AF have no definite cause, it may be the result of various other problems. Hence, kidney function and electrolytes are routinely determined, as well as thyroid-stimulating hormone (commonly suppressed in hyperthyroidism and of relevance if amiodarone is administered for treatment) and a blood count.
Electrocardiogram
Atrial fibrillation is diagnosed on an electrocardiogram (ECG), an investigation performed routinely whenever an irregular heartbeat is suspected. Characteristic findings are the absence of P waves, with disorganized electrical activity in their place, and irregular R–R intervals due to irregular conduction of impulses to the ventricles. At very fast heart rates, atrial fibrillation may look more regular, which may make it more difficult to separate from other supraventricular tachycardias or ventricular tachycardia.
QRS complexes should be narrow, signifying that they are initiated by normal conduction of atrial electrical activity through the intraventricular conduction system. Wide QRS complexes are worrisome for ventricular tachycardia, although, in cases where there is a disease of the conduction system, wide complexes may be present in A-fib with a rapid ventricular response.
If paroxysmal AF is suspected, but an ECG during an office visit shows only a regular rhythm, AF episodes may be detected and documented with the use of ambulatory Holter monitoring (e.g., for a day). If the episodes are too infrequent to be detected by Holter monitoring with reasonable probability, then the person can be monitored for longer periods (e.g., a month) with an ambulatory event monitor.
Echocardiography
In general, a non-invasive transthoracic echocardiogram (TTE) is performed in newly diagnosed AF, as well as if there is a major change in the person's clinical state. This ultrasound-based scan of the heart may help identify valvular heart disease (which may greatly increase the risk of stroke and alter recommendations for the appropriate type of anticoagulation), left and right atrial size (which predicts the likelihood that AF may become permanent), left ventricular size and function, peak right ventricular pressure (pulmonary hypertension), presence of left atrial thrombus (low sensitivity), presence of left ventricular hypertrophy and pericardial disease.
Significant enlargement of both the left and right atria is associated with long-standing atrial fibrillation and, if noted at the initial presentation of atrial fibrillation, suggests that the atrial fibrillation is likely to be of a longer duration than the individual's symptoms.
Chest X-ray
In general, a chest X-ray is performed only if a pulmonary cause of atrial fibrillation is suggested, or if other cardiac conditions are suspected (in particular congestive heart failure). This may reveal an underlying problem in the lungs or the blood vessels in the chest.
Transesophageal echocardiogram
A regular echocardiogram (transthoracic echocardiogram; TTE) has a low sensitivity for identifying blood clots in the heart. If this is suspected (e.g. when planning urgent electrical cardioversion), a transesophageal echocardiogram (TEE or TOE where British spelling is used) is preferred.
The TEE has much better visualization of the left atrial appendage than transthoracic echocardiography. This structure, located in the left atrium, is the place where a blood clot forms in more than 90% of cases in non-valvular (or non-rheumatic) atrial fibrillation. TEE has a high sensitivity for locating thrombi in this area and can also detect sluggish blood flow in this area that is suggestive of blood clot formation. If a blood clot is seen on TEE, then cardioversion is contraindicated due to the risk of stroke, and anticoagulation is recommended.
Ambulatory Holter monitoring
A Holter monitor is a wearable ambulatory heart monitor that continuously monitors the heart rate and heart rhythm for a short duration, typically 24 hours. In individuals with symptoms of significant shortness of breath with exertion or palpitations regularly, a Holter monitor may be of benefit to determine whether rapid heart rates (or unusually slow heart rates) during atrial fibrillation are the cause of the symptoms.
Exercise stress testing
Some individuals with atrial fibrillation do well with normal activity but develop shortness of breath with exertion. It may be unclear whether the shortness of breath is due to a blunted heart rate response to exertion caused by excessive atrioventricular node-blocking agents, a very rapid heart rate during exertion, or other underlying conditions such as chronic lung disease or coronary ischemia. An exercise stress test will evaluate the individual's heart rate response to exertion and determine whether the AV node blocking agents are contributing to the symptoms.
Non-contact detection
Non-contact AF detection technologies have been developed that utilize pulsatile changes in skin color using video photoplethysmography (VPG). These methods can allow for remote AF detection (e.g., via Telehealth), which can improve outcomes. Under optimal conditions — such as lighting (≥100 lux) and camera position — this method can achieve low error rates through computation of Pulse Harmonic Strength (PHS), a metric of pulse rate. However, due to limited initial studies, little is known about the potential computational-based racial biases embedded in these algorithms. Algorithmic biases have been found in facial detection software (Gender Shades); these technologies may be inequitable and exacerbate healthcare disparities. More clinical validation is needed before these tools become widely used for medical decisions.
Classification
The American College of Cardiology (ACC), American Heart Association (AHA), and the European Society of Cardiology (ESC) recommend in their guidelines the following classification system based on simplicity and clinical relevance.
All people with AF are initially in the category called first detected AF. These people may or may not have had previous undetected episodes. If a first detected episode stops on its own in less than seven days and then another episode begins, later on, the category changes to paroxysmal AF. Although people in this category have episodes lasting up to seven days, in most cases of paroxysmal AF, the episodes will stop in less than 24 hours. If the episode lasts for more than seven days, it is unlikely to stop on its own and is then known as persistent AF. In this case, cardioversion can be attempted to restore a normal rhythm. If an episode continues for a year or more, the rhythm is then known as longstanding persistent AF. If a decision is made by the person and their medical team to accept persistent AF and not attempt restoration of a normal sinus rhythm but instead manage the AF by simply controlling the person's ventricular rate then the rhythm is referred to as permanent AF. As a further subtype, AF that is detected only by an implanted or wearable cardiac monitor is known as subclinical AF.
Episodes that last less than 30 seconds are not considered in this classification system. Also, this system does not apply to cases where the AF is a secondary condition that occurs in the setting of a primary condition that may be the cause of the AF.
About half of people with AF have permanent AF, while a quarter have paroxysmal AF, and a quarter have persistent AF.
In addition to the above AF categories, which are mainly defined by episode timing and termination, the ACC/AHA/ESC guidelines describe additional AF categories in terms of other characteristics of the person. Valvular AF refers to AF attributable to moderate to severe mitral valve stenosis or atrial fibrillation in the presence of a mechanical artificial heart valve. This distinction may be useful as it has implications on appropriate treatment, including differing recommendations for anticoagulation, but some guidelines discourage use of this term as it may be confusing. Other historically used definitions include lone AF – AF occurring in those aged under 60 in the absence of other cardiovascular or respiratory diseases. This description is also discouraged as it is recognised that AF always has an underlying cause. Secondary AF refers to AF that occurs in the setting of another condition that have caused the AF, such as acute myocardial infarction, cardiac surgery, pericarditis, myocarditis, hyperthyroidism, pulmonary embolism, pneumonia, or another acute pulmonary disease.
Prevention
Prevention of atrial fibrillation focuses primarily on preventing or controlling its risk factors. Many of its risk factors, such as obesity, smoking, lack of physical activity, and excessive alcohol consumption, are modifiable and preventable with lifestyle modification or can be managed by a healthcare professional.
Lifestyle modification
Several healthy lifestyle behaviors are associated with a lower likelihood of developing atrial fibrillation. Accordingly, consensus guidelines recommend abstaining from alcohol and recreational drugs, stopping tobacco use, maintaining a healthy weight, and regularly participating in moderate-intensity physical activities. Consistent moderate-intensity aerobic exercise, defined as achieving 3.0–5.9 METs of intensity, for at least 150 minutes per week may reduce the risk of developing new-onset atrial fibrillation. Few studies have examined the role of specific dietary changes and how it relates to the prevention of atrial fibrillation.
Management
The main goals of treatment are to prevent circulatory instability and stroke. Rate or rhythm control is used to achieve the former, whereas anticoagulation is used to decrease the risk of the latter. If cardiovascularly unstable due to uncontrolled tachycardia, immediate cardioversion is indicated. Many antiarrhythmics, when used long term, increase the risk of death without any meaningful benefit. An integrated management approach, which includes stroke prevention, symptoms control and management of associated comorbidities has been associated with better outcomes in patients with atrial fibrillation.
This holistic or integrated care approach is summed up as the ABC (Atrial fibrillation Better Care) pathway, as follows:
A: Avoid stroke with Anticoagulation, where the default is stroke prevention unless the patient is at low risk. Stroke prevention means use of oral anticoagulation (OAC), whether with well managed vitamin K antagonists (VKA), with time in therapeutic range >70%, or more commonly, label-adherent dosed direct oral anticoagulant (DOAC).
B: Better symptom and atrial fibrillation management with patient-centred, symptom directed decisions on rate control or rhythm control. In some selected patients, use early rhythm control may be beneficial.
C: Cardiovascular risk factor and comorbidity management, including attention to lifestyle factors and psychological morbidity.
Lifestyle modification
Regular aerobic exercise improves atrial fibrillation symptoms and AF-related quality of life. The effect of high-intensity interval training on reducing atrial fibrillation burden is unclear. Weight loss of at least 10% is associated with reduced atrial fibrillation burden in people who are overweight or obese.
Comorbidity treatment
For people who have both atrial fibrillation and obstructive sleep apnea, observational studies suggest that continuous positive airway pressure (CPAP) treatment appears to lower the risk of atrial fibrillation recurrence after undergoing ablation. Randomized controlled trials examining the role of obstructive sleep apnea treatment on atrial fibrillation incidence and burden are lacking. Guideline-recommended lifestyle and medical interventions are recommended for people with atrial fibrillation and coexisting conditions such as hyperlipidemia, diabetes mellitus, or hypertension without specific blood sugar or blood pressure targets for people with atrial fibrillation.
Bariatric surgery may reduce the risk of new-onset atrial fibrillation in people with obesity without AF and may reduce the risk of a recurrence of AF after an ablation procedure in people with coexisting obesity and atrial fibrillation. It is important for all people with atrial fibrillation to optimize the control of all coexisting medical conditions that can worsen their atrial fibrillation, such as hyperthyroidism, diabetes, congestive heart failure, high blood pressure, chronic obstructive pulmonary disease, stimulant use (e.g., methamphetamine dependence), and excessive alcohol consumption.
Anticoagulants
Anticoagulation can be used to reduce the risk of stroke from AF. Anticoagulation is recommended in most people other than those at low risk of stroke or those at high risk of bleeding. The risk of falls and consequent bleeding in frail elderly people should not be considered a barrier to initiating or continuing anticoagulation since the risk of fall-related brain bleeding is low and the benefit of stroke prevention often outweighs the risk of bleeding. Similarly, the presence or absence of AF symptoms does not determine whether a person warrants anticoagulation and is not an indicator of stroke risk. Oral anticoagulation is underused in atrial fibrillation, while aspirin is overused in many who should be treated with a direct oral anticoagulant (DOAC) or warfarin. In 2019, DOACs were often recommended over warfarin by the American Heart Association.
The risk of stroke from non-valvular AF can be estimated using the CHA2DS2-VASc score. In the 2019 AHA/ACC/HRS guidelines, anticoagulation is recommended in non-valvular AF if there is a score of two or more in men and three or more in women and may be considered if there is a score of one in men or two in women; not using anticoagulation is reasonable if there is a score of zero in men or one in women. Guidelines from the American College of Chest Physicians, Asia-Pacific Heart Rhythm Society, Canadian Cardiovascular Society, European Society of Cardiology, Japanese Circulation Society, Korean Heart Rhythm Society, and the National Institute for Health and Care Excellence recommend the use of novel oral anticoagulants or warfarin with a CHA2DS2-VASc score of one over aspirin and some directly recommend against aspirin. Experts generally advocate for most people with atrial fibrillation with CHA2DS2-VASc scores of one or more receiving anticoagulation though aspirin is sometimes used for people with a score of one (moderate risk for stroke). There is little evidence to support the idea that the use of aspirin significantly reduces the risk of stroke in people with atrial fibrillation. Furthermore, aspirin's major bleeding risk (including bleeding in the brain) is similar to that of warfarin and DOACs despite its inferior efficacy.
Anticoagulation can be achieved through several means including warfarin, heparin, dabigatran, rivaroxaban, edoxaban, and apixaban. Many issues should be considered related to their comparative effectiveness, including the cost of DOACs, risk of stroke, risk of falls, comorbidities (such as chronic liver or kidney disease), the presence of significant mitral stenosis or mechanical heart valves, compliance, and speed of the desired onset of anticoagulation. The optimal approach to anticoagulation in people with AF and who simultaneously have other diseases (e.g., cirrhosis and end-stage kidney disease on dialysis) that predispose a person to both bleeding and clotting complications is unclear.
For those with non-valvular atrial fibrillation, DOACs are at least as effective as warfarin for preventing strokes and blood clots embolizing to the systemic circulation (if not more so) and are generally preferred over warfarin. DOACs carry a lower risk of bleeding in the brain compared to warfarin, although dabigatran is associated with a higher risk of intestinal bleeding. Dual antiplatelet therapy with aspirin and clopidogrel is inferior to warfarin for preventing strokes or systemic embolism and has comparable bleeding risk in people with atrial fibrillation. In those who are also on aspirin, however, DOACs appear to be better than warfarin.
Time in therapeutic range (TTR) and INR variability are commonly used to assess the quality of VKA treatment. Patients who are unable to maintain a therapeutic INR on VKA, as indicated by low TTR and/or high INR variability, are at an increased risk of thromboembolic and bleeding events. In these patients, treatment with a DOAC is recommended. While there are no significant changes in adherence, persistence or clinical outcomes in patients switched from a VKA to a DOAC, an increase in therapy satisfaction has been reported.
DOAC therapy is not recommended for all patients with atrial fibrillation. For instance, warfarin is the recommended anticoagulant for patients with atrial fibrillation who have mechanical heart valves.
Rate versus rhythm control
There are two ways to approach atrial fibrillation using medications: rate control and rhythm control. Both methods have similar outcomes. Rate control lowers the heart rate closer to normal, usually 60 to 100 bpm, without trying to convert to a regular rhythm. Rhythm control tries to restore a normal heart rhythm in a process called cardioversion and maintains the normal rhythm with medications. Studies suggest that rhythm control is more important in the acute setting AF, whereas rate control is more important in the chronic phase.
The risk of stroke appears to be lower with rate control versus attempted rhythm control, at least in those with heart failure. AF is associated with a reduced quality of life, and, while some studies indicate that rhythm control leads to a higher quality of life, some did not find a difference.
Neither rate nor rhythm control is superior in people with heart failure when they are compared in various clinical trials. However, rate control is recommended as the first-line treatment regimen for people with heart failure. On the other hand, rhythm control is only recommended when people experience persistent symptoms despite adequate rate control therapy.
In those with a fast ventricular response, intravenous magnesium significantly increases the chances of achieving successful rate and rhythm control in the urgent setting without major side-effects. A person with poor vital signs, mental status changes, preexcitation, or chest pain often will go to immediate treatment with synchronized DC cardioversion. Otherwise, the decision of rate control versus rhythm control using medications is made. This is based on several criteria that include whether or not symptoms persist with rate control.
Rate control
Rate control to a target heart rate of fewer than 110 beats per minute is recommended in most people. Lower heart rates may be recommended in those with left ventricular hypertrophy or reduced left ventricular function. Rate control is achieved with medications that work by increasing the degree of the block at the level of the AV node, decreasing the number of impulses that conduct into the ventricles. This can be done with:
Beta blockers (preferably the "cardioselective" beta blockers such as metoprolol, bisoprolol, or nebivolol)
Non-dihydropyridine calcium channel blockers (e.g., diltiazem or verapamil)
Cardiac glycosides (e.g., digoxin) – have less use, apart from in older people who are sedentary. They are not as effective as either beta-blockers or calcium channel blockers.
Patients with chronic AF are recommended to take either beta blockers or calcium channel blockers.
In addition to these agents, amiodarone has some AV node blocking effects (in particular when administered intravenously) and can be used in individuals when other agents are contraindicated or ineffective (particularly due to hypotension).
Cardioversion
Cardioversion is the attempt to switch an irregular heartbeat to a normal heartbeat using electrical or chemical means.
Electrical cardioversion involves the restoration of normal heart rhythm through the application of a DC electrical shock. The exact placement of the pads does not appear to be important.
Chemical cardioversion is performed with medications, such as amiodarone, dronedarone, procainamide (especially in pre-excited atrial fibrillation), dofetilide, ibutilide, propafenone, or flecainide.
After successful cardioversion, the heart may be stunned, which means that there is a normal rhythm, but the restoration of normal atrial contraction has not yet occurred.
Surgery
Ablation
Catheter ablation (CA) is a procedure performed by an electrophysiologist, a cardiologist who specializes in heart rhythm problems, to restore the heart's normal rhythm by destroying, or electrically isolating, specific parts of the atria. A group of cardiologists led by Dr Haïssaguerre from Bordeaux University Hospital noted in 1998 that the pulmonary veins are an important source of ectopic beats, initiating frequent paroxysms of atrial fibrillation, with these foci responding to treatment with radio-frequency ablation. Most commonly, CA electrically isolates the left atrium from the pulmonary veins, where most of the abnormal electrical activity promoting atrial fibrillation originates. CA is a form of rhythm control that restores normal sinus rhythm and reduces AF-associated symptoms more reliably than antiarrhythmic medications.
Electrophysiologists generally use two forms of catheter ablation—radiofrequency ablation, or cryoablation. In young people with little-to-no structural heart disease where rhythm control is desired and cannot be maintained by medication or cardioversion, radiofrequency catheter ablation or cryoablation may be attempted and may be preferred over several years of medical therapy. Although radiofrequency ablation has become an accepted intervention in selected younger people and may be more effective than medication at improving symptoms and quality of life, there is no evidence that ablation reduces all-cause mortality, stroke, or heart failure. Some evidence indicates CA may be particularly helpful for people with AF who also have heart failure. AF may recur in people who have undergone CA and nearly half of people who undergo it will require a repeat procedure to achieve long-term control of their AF.
In general, CA is more successful at preventing AF recurrence if AF is paroxysmal as opposed to persistent. As CA does not reduce the risk of stroke, many are advised to continue their anticoagulation. Possible complications include common, minor complications such as the formation of a collection of blood at the site where the catheter goes into the vein (access site hematoma), but also more dangerous complications including bleeding around the heart (cardiac tamponade), stroke, damage to the esophagus (atrio-esophageal fistula), or even death. Use of pulsed field ablation as a non-thermal method of inducing electroporation avoids damage to the phrenic nerve, esophagus, and blood vessels, while being at least as effective as thermal ablation methods.
A hybrid convergent procedure has been developed which combines endocardial ablation with epicardial ablation, which can reduce AF recurrence to less than 5% for over one year. The epicardial ablation is performed first, with a minimally invasive surgical approach.
Maze procedure
An alternative to catheter ablation is surgical ablation. The maze procedure, first performed in 1987, is an effective invasive surgical treatment that is designed to create electrical blocks or barriers in the atria of the heart. The idea is to force abnormal electrical signals to move along one, uniform path to the lower chambers of the heart (ventricles), thus restoring the normal heart rhythm. People with AF often undergo cardiac surgery for other underlying reasons and are frequently offered concomitant AF surgery to reduce the frequency of short- and long-term AF. Concomitant AF surgery is more likely to lead to the person being free from atrial fibrillation and off medications long-term after surgery and Cox-Maze IV procedure is the gold standard treatment. There is a slightly increased risk of needing a pacemaker following the procedure. Less invasive modifications of the maze procedure have been developed, designated as minimaze procedures.
Left atrial appendage occlusion
There is growing evidence that left atrial appendage occlusion therapy may reduce the risk of stroke in people with non-valvular AF as much as warfarin. The addition of left atrial appendage isolation to catheter ablation has reduced AF recurrence by 80% in patients with persistent AF.
After surgery
After catheter ablation, people are moved to a cardiac recovery unit, intensive care unit, or cardiovascular intensive care unit where they are not allowed to move for 46 hours. Minimizing movement helps prevent bleeding from the site of the catheter insertion. The length of time people stay in the hospital varies from hours to days. This depends on the problem, the length of the operation, and whether or not general anesthetic was used. Additionally, people should not engage in strenuous physical activityto maintain a low heart rate and low blood pressurefor around six weeks.
AF often occurs after cardiac surgery and is usually self-limiting. It is strongly associated with age, preoperative hypertension, and the number of vessels grafted. Measures should be taken to control hypertension preoperatively to reduce the risk of AF. Also, people with a higher risk of AF, e.g., people with pre-operative hypertension, more than three vessels grafted, or greater than 70 years of age, should be considered for prophylactic treatment. Postoperative pericardial effusion is also suspected to be the cause of atrial fibrillation. Prophylaxis may include prophylactic postoperative rate and rhythm management. Some authors perform posterior pericardiotomy to reduce the incidence of postoperative AF. When AF occurs, management should primarily be rate and rhythm control. However, cardioversion may be used if the patient is hemodynamically unstable, highly symptomatic, or AF persists for six weeks after discharge. In persistent cases, anticoagulation should be used.
Prognosis
Atrial fibrillation can progress from infrequent occurrences to more frequent occurrences, ultimately becoming permanent. Some cases do not progress, especially among patients with a healthy lifestyle.
Many mechanisms contribute to cardiac remodeling leading to a worsening of atrial fibrillation, including fibrosis, fatty infiltration, amyloidosis, and ion channel modifications. Fatty infiltration helps explain why obesity is a risk factor for atrial fibrillation in one fifth of patients.
Atrial fibrillation increases the risk of heart failure by 11 per 1000, kidney problems by 6 per 1000, death by 4 per 1000, stroke by 3 per 1000, and coronary heart disease by 1 per 1000. Women have a worse outcome overall than men. Evidence increasingly suggests that atrial fibrillation is independently associated with a higher risk of developing dementia.
Blood clots
Prediction of embolism
Among Danish men aged 50, with no risk factors, the 5-year risk of stroke was 1.1% and with AF alone 2.5%. For women the risks were slightly less, 0.7% and 2.1%. For men aged 70, the 5-year risk of stroke was 4.8% and with AF alone 6.8%. For women aged 70 the risk was again lower than for men, 3.4% with no added risk factor and 8.2% with AF.
Determining the risk of an embolism causing a stroke is important for guiding the use of anticoagulants. The most accurate clinical prediction rules are:
CHADS2
CHA2DS2-VASc score
Both the CHADS2 and the CHA2DS2-VASc score predict future stroke risk in people with A-fib with CHA2DS2-VASc score being more accurate. The addition of blood based biomarkers such as NT-proBNP and neurofilament light chain improves risk prediction significantly. Some that had a CHADS2 score of zero had a CHA2DS2-VASc score of three, with a 3.2% annual risk of stroke. Thus, a CHA2DS2-VASc score of zero is considered very low risk.
Mechanism of thrombus formation
In atrial fibrillation, the lack of an organized atrial contraction can result in some stagnant blood in the left atrium (LA) or left atrial appendage (LAA). This lack of movement of blood can lead to thrombus formation (blood clotting). If the clot becomes mobile and is carried away by the blood circulation, it is called an embolus. An embolus proceeds through smaller and smaller arteries until it plugs one of them and prevents blood from flowing through the artery. This process results in end organ damage due to the loss of nutrients, oxygen, and the removal of cellular waste products. Emboli in the brain may result in an ischemic stroke or a transient ischemic attack (TIA).
More than 90% of cases of thrombi associated with non-valvular atrial fibrillation evolve in the left atrial appendage. However, the LAA lies in close relation to the free wall of the left ventricle, and thus the LAA's emptying and filling, which determines its degree of blood stagnation, may be helped by the motion of the wall of the left ventricle if there is good ventricular function.
Dementia
Atrial fibrillation has been independently associated with a higher risk of developing cognitive impairment, vascular dementia, and Alzheimer disease and with elevated levels of neurofilament light chain in blood, a biomarker indicating neuroaxonal injury. Several mechanisms for this association have been proposed, including silent small blood clots (subclinical microthrombi) traveling to the brain resulting in small ischemic strokes without symptoms, altered blood flow to the brain, inflammation, clinically silent small bleeds in the brain, and genetic factors. Tentative evidence suggests that effective anticoagulation with direct oral anticoagulants or warfarin may be somewhat protective against AF-associated dementia and evidence of silent ischemic strokes on MRI but this remains an active area of investigation.
Epidemiology
Atrial fibrillation is the most common arrhythmia and affects more than 33 million people worldwide. In Europe and North America, , it affects about 2% to 3% of the population. This is an increase from 0.4 to 1% of the population around 2005. In the developing world, rates are about 0.6% for males and 0.4% for females. The number of people diagnosed with AF has increased due to better detection of silent AF and increasing age and conditions that predispose to it.
It also accounts for one-third of hospital admissions for cardiac rhythm disturbances, and the rate of admissions for AF has risen. AF is the cause for 20% to 30% of all ischemic strokes. After a transient ischemic attack or stroke, about 11% are found to have a new diagnosis of atrial fibrillation. 3% to 11% of patients with AF have structurally normal hearts. Approximately 2.2 million individuals in the United States and 4.5 million in the European Union have AF.
The number of new cases each year of AF increases with age. In people older than 80 years, it affects about 8%. In contrast, in younger people the prevalence is estimated to be 0.05% and is associated with congenital heart disease or structural heart disease in this demographic. As of 2001, it was anticipated that in developed countries, the number of people with atrial fibrillation was likely to increase during the following 50 years, due to the growing proportion of elderly people.
Gender
Atrial fibrillation is more common in men than in women when reviewed in European and North American populations. In developed and developing countries, there is also a higher rate in men than in women. The risk factors associated with AF are also distributed differently according to gender. In men, coronary disease is more frequent, while in women, high systolic blood pressure and valvular heart disease are more prevalent.
Ethnicity
Rates of AF are lower in populations of African descent than in populations of European descent. African descent is associated with a protective effect for AF, due to the lower presence of SNPs with guanine alleles. European ancestry has more frequent mutations. The variant rs4611994 for the gene PITX2 is associated with risk of AF in African and European populations. Hispanic and Asian populations have a lower risk of AF than European populations. The risk of AF in non-European populations is associated with characteristic risk factors of these populations, such as hypertension.
Young people
Atrial fibrillation is an uncommon condition in children but sometimes occurs in association with certain inherited and acquired conditions. Congenital heart disease and rheumatic fever are the most common causes of atrial fibrillation in children. Other inherited heart conditions associated with the development of atrial fibrillation in children include Brugada syndrome, short QT syndrome, Wolff Parkinson White syndrome, and other forms of supraventricular tachycardia (e.g., AV nodal reentrant tachycardia). Adults who survived congenital heart disease have an increased risk of developing AF. In particular, people who had atrial septal defects, Tetralogy of Fallot, or Ebstein's anomaly, and those who underwent the Fontan procedure, are at higher risk with prevalence rates of up to 30% depending on the heart's anatomy and the person's age.
History
Because the diagnosis of atrial fibrillation requires measurement of the electrical activity of the heart, atrial fibrillation was not truly described until 1874, when Edmé Félix Alfred Vulpian observed the irregular atrial electrical behavior that he termed "fremissement fibrillaire" in dog hearts. In the mid-18th century, Jean-Baptiste de Sénac made note of dilated, irritated atria in people with mitral stenosis. The irregular pulse associated with AF was first recorded in 1876 by Carl Wilhelm Hermann Nothnagel and termed "delirium cordis", stating that "[I]n this form of arrhythmia the heartbeats follow each other in complete irregularity. At the same time, the height and tension of the individual pulse waves are continuously changing". Correlation of delirium cordis with the loss of atrial contraction, as reflected in the loss of a waves in the jugular venous pulse, was made by Sir James MacKenzie in 1904. Willem Einthoven published the first ECG showing AF in 1906. The connection between the anatomic and electrical manifestations of AF and the irregular pulse of delirium cordis was made in 1909 by Carl Julius Rothberger, Heinrich Winterberg, and Sir Thomas Lewis.
Other animals
Atrial fibrillation occurs in other animals, including cats, dogs, and horses. Unlike humans, dogs rarely develop the complications that stem from blood clots breaking off from inside the heart and traveling through the arteries to distant sites (thromboembolic complications). Cats rarely develop atrial fibrillation but appear to have a higher risk of thromboembolic complications than dogs.
Cats and dogs with atrial fibrillation often have underlying structural heart disease that predisposes them to the condition. The medications used in animals for atrial fibrillation are largely similar to those used in humans. Electrical cardioversion is occasionally performed in these animals, but the need for general anesthesia limits its use. Standardbred horses appear to be genetically susceptible to developing atrial fibrillation. Horses that develop atrial fibrillation often have minimal or no underlying heart disease, and the presence of atrial fibrillation in horses can adversely affect physical performance.
| Biology and health sciences | Cardiovascular disease | Health |
1875639 | https://en.wikipedia.org/wiki/Dashpot | Dashpot | A dashpot, also known as a damper, is a mechanical device that resists motion via viscous friction. The resulting force is proportional to the velocity, but acts in the opposite direction, slowing the motion and absorbing energy. It is commonly used in conjunction with a spring. The process and instrumentation diagram (P&ID) symbol for a dashpot is .
Types
The two most common types of dashpots are linear and rotary.
Linear damper
Linear dashpots — or linear dampers — are used to exert a force opposite to a translation movement. They are generally specified by stroke (amount of linear displacement) and damping coefficient (force per velocity).
Rotary damper
Similarly, rotary dampers will tend to oppose any torque applied to them, in an amount proportional to their rotational speed. Their damping coefficients will usually be specified by torque per angular velocity. One can distinguish two kinds of viscous rotary dashpots:
Vane dashpots which have a limited angular range but provide a significant damping torque. The damping force is the result of one or multiple vanes moving through a viscous fluid and letting it flow via calibrated openings.
Continuous rotation dashpots which aren't limited in their rotation angle but provide a smaller damping coefficient. These use the friction generated by the shearing forces induced in the viscous fluid itself by the difference in motion between the dashpot's rotor and stator.
Eddy current damper
A less common type of dashpot is an eddy current damper, which uses a large magnet inside a tube constructed of a non-magnetic but conducting material (such as aluminium or copper). Like a common viscous damper, the eddy current damper produces a resistive force proportional to velocity. A common use of the eddy current damper is in balance scales. This is a frictionless method that allows the scale to quickly come to rest.
One-way operation
Dashpots frequently use a one-way mechanical bypass to permit fast unrestricted motion in one direction and slow motion using the dashpot in the opposite direction. This permits, for example, a door to be opened quickly without added resistance, but then to close slowly using the dashpot. For hydraulic dashpots this unrestricted motion is accomplished using a one-way check-valve that allows fluid to bypass the dashpot fluid constriction. Non-hydraulic rotatory dashpots may use a ratcheting gear to permit free motion in one direction.
Theory
Dashpots are frequently used to add damping to dynamic systems. When designing and analyzing systems, dashpots are often assumed to be linear, meaning that their output force is linearly proportional to their velocity.
This permits convenient analysis of systems such as harmonic oscillators. However, the behavior of real-world dashpots is frequently non-linear, meaning that the force is proportional to the velocity raised to some exponent α, which can vary between 0.2 and 2.
Different exponents are better suited for different applications, but exponents other than 1.0 must be analyzed with numeric methods instead of calculus.
Applications
A dashpot is a common component in a door closer to prevent it from slamming shut. A spring applies force to close the door, which the dashpot offsets by forcing fluid to flow through an orifice, often adjustable, between reservoirs, which slows the motion of the door.
Consumer electronics often use dashpots where it is undesirable for a media access door or control panel to suddenly pop open when the door latch is released. The dashpot provides a steady, gentle motion until the access door has fully opened.
Dashpots are commonly used in dampers and shock absorbers. The hydraulic cylinder in an automobile shock absorber is a dashpot.
They are also used on carburetors, where the return of the throttle lever is cushioned just before the throttle fully closes, then is allowed to fully close slowly to reduce emissions. The British SU carburettor's main piston carries a stepped needle. This needle is held in the fuel flow orifice. The manifold vacuum causes this piston to rise allowing more fuel into the airflow. The SU's dashpot has a fixed hydraulic piston, damping the main piston as it moves upward. A valve in the piston disables the damping as the main piston returns.
Large forces and high speeds can be controlled by dashpots. For example, they are used to arrest the steam catapults on aircraft carrier decks.
Relays can be made to have a long delay by utilizing a piston filled with fluid that is allowed to escape slowly. Electrical switchgear may use dashpots in their overcurrent sensing mechanism to reduce reaction speed to brief events, thus making them less sensitive to false-triggering during transients whilst still remaining sensitive to sustained overloads. Another use is for delaying the closing or opening of an electrical circuit. Such a dashpot timer might be used for example for timed staircase lighting.
Anti-stall mechanisms in internal combustion engines are aimed to prevent stalling of the engine at low rpm. Anti-stall mechanisms use dashpots to arrest the final closing movement of the throttle.
Large dashpots are added to bridges and buildings to protect against earthquakes and wind vibrations.
Viscoelasticity
Dashpots are used as models of materials that exhibit viscoelastic behavior, such as muscle tissue. Maxwell and Kelvin–Voigt models of viscoelasticity use springs and dashpots in series and parallel circuits respectively. Models containing dashpots add a viscous, time-dependent element to the behavior of solids, allowing complex behaviors like creep and stress relaxation to be modeled.
| Technology | Mechanisms | null |
1877567 | https://en.wikipedia.org/wiki/Non-nucleophilic%20base | Non-nucleophilic base | As the name suggests, a non-nucleophilic base is a sterically hindered organic base that is a poor nucleophile. Normal bases are also nucleophiles, but often chemists seek the proton-removing ability of a base without any other functions. Typical non-nucleophilic bases are bulky, such that protons can attach to the basic center but alkylation and complexation is inhibited.
Non-nucleophilic bases
A variety of amines and nitrogen heterocycles are useful bases of moderate strength (pKa of conjugate acid around 10-13)
N,N-Diisopropylethylamine (DIPEA, also called Hünig's Base), pKa = 10.75
1,8-Diazabicycloundec-7-ene (DBU) - useful for E2 elimination reactions, pKa = 13.5
1,5-Diazabicyclo(4.3.0)non-5-ene (DBN) - comparable to DBU
2,6-Di-tert-butylpyridine, a weak non-nucleophilic base pKa = 3.58
Phosphazene bases, such as t-Bu-P4
Non-nucleophilic bases of high strength are usually anions. For these species, the pKas of the conjugate acids are around 35–40.
Lithium diisopropylamide (LDA), pKa = 36
Silicon-based amides, such as sodium and potassium bis(trimethylsilyl)amide (NaHMDS and KHMDS, respectively)
Lithium tetramethylpiperidide (LiTMP or harpoon base)
Other strong non-nucleophilic bases are sodium hydride and potassium hydride. These compounds are dense, salt-like materials that are insoluble and operate by surface reactions.
Some reagents are of high basicity (pKa of conjugate acid around 17) but of modest but not negligible nucleophilicity. Examples include sodium tert-butoxide and potassium tert-butoxide.
Example
The following diagram shows how the hindered base, lithium diisopropylamide, is used to deprotonate an ester to give the enolate in the Claisen ester condensation, instead of undergoing a nucleophilic substitution.
This reaction (deprotonation with LDA) is commonly used to generate enolates.
| Physical sciences | Concepts | Chemistry |
1878290 | https://en.wikipedia.org/wiki/Potassium%20superoxide | Potassium superoxide | Potassium superoxide is an inorganic compound with the formula . It is a yellow paramagnetic solid that decomposes in moist air. It is a rare example of a stable salt of the superoxide anion. It is used as a scrubber, dehumidifier, and generator in rebreathers, spacecraft, submarines, and spacesuits.
Production and reactions
Potassium superoxide is produced by burning molten potassium in an atmosphere of excess oxygen.
The salt consists of and ions, linked by ionic bonding. The O–O distance is 1.28 Å.
Reactivity
Potassium superoxide is a source of superoxide, which is an oxidant and a nucleophile, depending on its reaction partner.
Upon contact with water, it undergoes disproportionation to potassium hydroxide, oxygen, and hydrogen peroxide:
It reacts with carbon dioxide, releasing oxygen:
Theoretically, 1 kg of absorbs 0.310 kg of while releasing 0.338 kg of . One mole of absorbs 0.5 moles of and releases 0.75 moles of oxygen.
Potassium superoxide finds only niche uses as a laboratory reagent. Because it reacts with water, is often studied in organic solvents. Since the salt is poorly soluble in nonpolar solvents, crown ethers are typically used. The tetraethylammonium salt is also known. Representative reactions of these salts involve using superoxide as a nucleophile, e.g., in converting alkyl bromides to alcohols and acyl chlorides to diacyl peroxides.
Ion exchange with tetramethylammonium hydroxide gives tetramethylammonium superoxide, a yellow solid.
Applications
The Russian Space Agency has successfully used potassium superoxide in chemical oxygen generators for its spacesuits and Soyuz spacecraft. Potassium superoxide was also used in a rudimentary life support system for five mice as part of the Biological Cosmic Ray Experiment on Apollo 17.
has also been used in canisters for rebreathers for firefighting and mine rescue, and in cartridges for chemical oxygen generators on submarines. A flash fire caused by dropping such a cartridge into seawater contributed to the Kursk disaster. This highly exothermic reaction with water is also the reason why potassium superoxide has had limited use in scuba rebreathers.
| Physical sciences | Oxide salts | Chemistry |
1879453 | https://en.wikipedia.org/wiki/Mercury%20sulfide | Mercury sulfide | Mercury sulfide, or mercury(II) sulfide is a chemical compound composed of the chemical elements mercury and sulfur. It is represented by the chemical formula HgS. It is virtually insoluble in water.
Crystal structure
HgS is dimorphic with two crystal forms:
red cinnabar (α-HgS, trigonal, hP6, P3221) is the form in which mercury is most commonly found in nature. Cinnabar has rhombohedral crystal system. Crystals of red are optically active. This is caused by the Hg-S helices in the structure.
black metacinnabar (β-HgS) is less common in nature and adopts the zinc blende crystal structure (T2d-F3m).
Preparation and chemistry
β-HgS precipitates as a black solid when Hg(II) salts are treated with H2S. The reaction is conveniently conducted with an acetic acid solution of mercury(II) acetate. With gentle heating of the slurry, the black polymorph converts to the red form. β-HgS is unreactive to all but concentrated acids.
Mercury is produced from the cinnabar ore by roasting in air and condensing the vapour.
HgS → Hg + S
Uses
When α-HgS is used as a red pigment, it is known as cinnabar. The tendency of cinnabar to darken has been ascribed to conversion from red α-HgS to black β-HgS. However β-HgS was not detected at excavations in Pompeii, where originally red walls darkened, and was attributed to the formation of Hg-Cl compounds (e.g., corderoite, calomel, and terlinguaite) and calcium sulfate, gypsum.
As the mercury cell as used in the chlor-alkali industry (Castner–Kellner process) is being phased out over concerns over mercury emissions, the metallic mercury from these setups is converted into mercury sulfide for underground storage.
With a band gap of 2.1 eV and its stability, it is possible to be used as photoelectrochemical cell.
Neutralization with sulfur has been suggested to clean mercury spills, but the reaction does not proceed rapidly and completely enough for emergencies.
| Physical sciences | Sulfide salts | Chemistry |
1880422 | https://en.wikipedia.org/wiki/Oncorhynchus | Oncorhynchus | Oncorhynchus is a genus of ray-finned fish in the subfamily Salmoninae of the family Salmonidae, native to coldwater tributaries of the North Pacific basin. The genus contains twelve extant species, namely six species of Pacific salmon and six species of Pacific trout, all of which are migratory (either anadromous or potamodromous) mid-level predatory fish that display natal homing and semelparity.
The name of the genus is derived from the Greek (, 'lump, bend') + (, 'snout'), in reference to the hooked secondary sexual characteristic — known as the kype — that the males develop on the lower jaw tip during mating season.
Range
Salmon and trout within Onchorhynchus are native to the tributaries of the North Pacific Ocean, with their native ranges extend from the Bering Sea coasts southwards to as far as Taiwan in the west and Mexico in the east, although most of them are distributed in high-latitude cold waters from the Russian Far East to the Pacific Northwest. In North America, some subspecies of cutthroat trout (O. clarkii) have become landlocked populations native to endorheic waters in the Rocky Mountains and the Great Basin, while others have crossed the Continental Divide to inhabit the Rio Grande and western tributaries of the Mississippi River, both of which drain into the Gulf of Mexico instead of the Pacific Ocean. Several species of Oncorhynchus, such as the rainbow trout (O. mykiss) and Chinook salmon (O. tshawytscha), have been widely introduced into non-native waters around the globe, establishing self-sustaining wild populations.
The six Pacific salmons of Oncorhynchus are anadromous (migratory) and semelparous (die after spawning). Migration can be affected by parasites. Infected individuals can become weak and probably have shortened lifespans. Infection with parasites creates an effect known as culling whereby fish that are infected are less likely to complete the migration. Anadromous forms of Oncorhynchus mykiss known as steelhead are iteroparous. The coastal cutthroat trout (Oncorhynchus clarkii clarkii) is considered semi-anadromous, as it spends some time in the ocean, usually much closer to its native stream than its fully anadromous relatives.
Evolution
Several Late Miocene (about 7 million years ago (Mya)) trout-like fossils in Idaho, in the Clarkia Lake beds, appear to be of Oncorhynchus. The presence of these species so far inland established Oncorhynchus was not only present in the Pacific drainages before the beginning of the Pliocene (5–6 Mya), but also that rainbow and cutthroat trout, and Pacific salmon lineages had diverged before the beginning of the Pliocene. Consequently, the split between Oncorhynchus and Salmo (Atlantic salmon) must have occurred well before the Pliocene. Suggested dates have gone back as far as the Early Miocene (about 20 Mya). One fossil species assigned to this genus, O. rastrosus, the spike-toothed salmon (synonym Smilodonichthys), is a long species known from Late Miocene to Pliocene fossils.
Speciation among Oncorhynchus has been examined for decades, and a family "tree" is not yet completely developed for the Pacific salmonids. Mitochondrial DNA (mtDNA) research has been completed on a variety of Pacific trout and salmonid species, but the results do not necessarily agree with fossil research, or molecular research. Chum, pink and sockeye salmon lineages are generally agreed to have diverged in the sequence after other species. Montgomery (2000) discusses the pattern of the fossil record as compared to tectonic shifts in the plates of the Pacific Northwest of America. The (potential) divergence in Oncorhynchus lineages appear to follow the uprising of the Pacific Rim. The climatic and habitat changes that would follow such a geologic event are discussed, in the context of potential stressors leading to adaptation and speciation.
One interesting case involving speciation with salmon is that of the kokanee salmon (landlocked sockeye salmon). Kokanee sockeye evolve differently from anadromous sockeye—they reach the level of "biological species". Biological species—as opposed to morphological species—are defined by the capacity to maintain themselves in sympatry as independent genetic entities. This definition can be vexing because it apparently applies only to sympatry, and this limitation makes the definition difficult to apply. Examples in Washington State, Canada, and elsewhere have two populations living in the same lake, but spawning in different substrates at different times, and eating different food sources. There is no pressure to compete or interbreed (two responses when resources are short). These types of kokanee salmon show the principal attributes of a biological species: they are reproductively isolated and show strong resources partitioning.
Decline of Oncorhynchus populations
A general decline in overall Pacific salmon populations began in the mid-19th century. As the result of western expansion and development in the U.S., experts estimate salmon populations in the Columbia River basin had been reduced to less than 20% of their pre-1850 levels by 1933. In 2008, Lackey estimated that Pacific salmon stocks in the Pacific Northwest were less than 10% of their pre-1850 numbers. Many of the remaining salmon runs are dominated by hatchery-raised salmon, not wild salmon. Many isolated subspecies of the Pacific trouts, particularly those of Oncorhynchus mykiss rainbow trout and
Oncorhynchus clarki cutthroat trout have declined in their native ranges. Many local populations or distinct population segments of anadromous forms of steelhead have declined in their native ranges. The resulting declines have resulted in a number of populations of Oncorhynchus species or subspecies being listed as either endangered, threatened or as "Species of Special Concern" by state, federal or international authorities. Two Oncorhynchus clarki subspecies are considered extinct. Declines are attributed to a wide variety of causes—overfishing, habitat loss and degradation, artificial propagation, stocking, and hybridization with or competition with introduced, non-native species. For example, the yellowfin cutthroat trout (Oncorhynchus clarki macdonaldi) is extinct as a result of the introduction of non-native rainbow trout into its native waters.
Pacific salmon are facing a widespread decline in body size. The mean body mass of sockeye salmon (O. nerka) decreased by 10% between 2000 and 2010. The mean body length of Oncorhynchus species decreased by 2–8% between 1990 and 2010. Salmon body size is decreasing due to a variety of evolutionary forces, including dams, fishing practices, climate change, and increased competition in the ocean. This trend in salmon size is expected to decrease nutrient cycling and salmon reproductive success while hurting the success of commercial fisheries and rural communities who rely on salmon for survival.
Influence of hatcheries
Declines in the abundance of wild salmon due to over fishing placed greater pressure on hatcheries to increase production and restore the wild salmon stock to supply fisheries. The problem is that hatcheries can never truly replicate the environment of wild salmon, an issue which often results in physiological and behavioral differences between wild salmon and those reared in hatcheries. These differences are often the product of genetic changes associated with inbreeding, artificial selection, and natural selection, as well as different environmental pressures acting on hatchery fish than wild populations. Due to the size selective nature of fishing favoring larger fish, a reduction in average size of the adult salmon has been observed over time. The smaller salmon make a greater proportion of the remaining individuals continuing the population, and problems arise when these hatchery-reared fish are introduced into the wild populations. Unlike wild salmon, larger salmon are selected for in hatcheries and are typically much larger than wild salmon. The result is that hatchery-produced salmon tend to out-compete wild salmon for space, food, and other resources. Some salmon species in hatcheries exhibit predatory behavior toward wild salmon because they grow to be much larger. Regardless of whether predation is observed, natural social interactions are disturbed by the release of large numbers of hatchery-reared salmon where wild populations are low because salmon in hatcheries naturally have a higher propensity towards aggressive behavior.
Overall, natural salmon populations are put at risk when hatchery-reared salmon populations are introduced due to competition for resources, predation by larger individuals, and negative social interactions that upset the natural order observed in wild salmon populations. As a result, wild salmon populations are steadily dropping as the pressure to continue breeding salmon in hatcheries increases. Conservation efforts that work to place limitations on hatcheries to increase the wild salmon populations are hindered by financial pressures because hatcheries effectively support many states economically by accounting for over 70% of the salmon harvested for recreational and commercial purposes.
Influence of overfishing
Pacific salmon are harvested throughout the world as a source of food in countries ranging from the United States to South Korea. Over the past century, Pacific salmon have been extensively fished through both recreational, artisan and commercial fishing. In fact, since the 1970s there has been a nearly threefold increase in catch of Pacific salmon. As this catch has increased, a selection of reduced body size has been observed. In Japanese chum salmon, for example, between 1970 and 1988 there has been a continuous decrease in body weight averaging between 11 and 32 percent. In part, this decline in body weight has been related to the size selective effect of fishing gear used in the harvesting of salmon populations. Salmon of larger body weight are more apt to be caught during fishing efforts, causing lower body weight to be a beneficial character trait for survival. Thus, Pacific salmon have become continuously smaller in body size. However, studies have also shown that for Pacific salmon, a larger mean size at the time of reproduction increases the survival of offspring. The life history of salmon favors delayed reproduction because fecundity increased with body size. Consequently, the smaller body size of salmon results in a negative impact on population growth by decreasing the survivability of progeny, and thus decreasing the growth rate of populations. This reduction of productivity in Pacific salmon is, in part, seeded in overfishing and has caused a reduction in population sizes throughout Pacific salmon species. Today, it seems that population numbers of Pacific salmon are on the rise; however, the consequences from the overfishing in the 70s and 80s are still being reflected, with the average body size of salmon being smaller than before the event of overfishing.
Conservation
Canadian efforts
There has been evidence that the sockeye salmon are affected by thermal conditions and their responses to temperature are relatively strong and tend to vary from region to region. Canada has also used the Species at Risk Act to recognize the importance of biological diversity when it comes to the conservation of the salmon population. This means that multiple species of salmon would be looked at when it comes to conservation as well as multiple areas that each species live in. COSEWIC, a Canadian organization for the conservation of species, has named the Interior Fraser River Coho, the Cultus Lake Sockeye, and the Sakinaw Lake Sockeye to all be endangered. In British Columbia sockeye salmon in four different watersheds were certified by the Marine Stewardship Council, or MSC, as sustainable fisheries in July 2010 and the certification is good for a period of five years. In 2011 MSC also certified the Pink Salmon Fishery and as of 2012 The Chum Salmon Fisheries started their review under the MSC to become certified as a sustainable fishery.
American efforts
The US government has been working to develop a nationwide policy for the salmon populations. The Pacific Salmon Stronghold Conservation Act was re-submitted to congress and if passed will create geographic strongholds for salmon populations. Other policies include the Wild Salmon Policy which was enacted in 2005; its number one focus is the conservation of salmon off of the coasts. Even localized policies have begun, with one in Oregon which focuses on the southernmost watershed and was approved January 2013. In the Alaskan efforts, there is evidence of eight known regional groups of survival. It is also seen that the emigration of smolts (young salmon) from freshwater to other areas such as marine areas have shown significant consequences on the survival of different salmon groups. The Alaska Department of Fish and Game first received MSC, Marine Stewardship Council, Certification in sustainable seafood back in 2000. Each certification is good for a period of five years, with yearly check ups to ensure that the fishery remains sustainable. It was renewed again in 2007, but in 2012 The ADFG left the program. The Annette Island Reserve salmon fishery is under the control of the Metlakatla Indian community and as such was not included in the previous assessments of the Alaskan fisheries. It received its sustainability certification in June 2011. The Wild Salmon Center is a nonprofit organization that works on promoting conservation efforts for salmon worldwide and in the United States; it has helped secure protected watershed areas for Russian and west coast salmon. Other efforts of the Wild Salmon Center include combating illegal fishing, maintaining sustainable fisheries, and creating local watersheds as new habitats.
Russian efforts
Poaching is a threat to Oncorhynchus salmon and steelhead populations in Russia. It is estimated that illegal catching of salmon is 1.5 times more than the reported catch. The Wild Salmon Center is working with Russian authorities to try to help improve traceability systems so that markets can distinguish between legal sustainable salmon and the illegal salmon. The Wild Salmon Center has secured some of its protected locations for the salmon populations. In efforts with the WWF, the Wild Salmon Center was also able to have a Sockeye Salmon fishery certified as completely sustainable in 2012. The Iturup Island Pink and Chum Salmon Fishery was first certified in 2009 and was the first Russian salmon fishery to receive certification in sustainability by MSC. Other fisheries that were certified by MSC include the Northeast Sakhalin Island Pink Salmon, certified in June 2012, and the Ozernaya River Sockeye Salmon, certified in September 2012. The Aniva Bay Pink Salmon and the Sakhalin Island Pink salmon are both under review by the MSC.
Introductions and aquaculture
Several species of Oncorhynchus have been successfully introduced into non-native waters, establishing self-sustaining wild populations. The Rainbow Trout Oncorhynchus mykiss is the most widely introduced species of the genus. Rainbow Trout, Chinook Salmon Oncorhynchus tshawytscha and Coho Salmon Oncorhynchus kisutch have established wild, self-sustaining populations in the Great Lakes and Chinook in New Zealand (known there as quinnat, king or spring salmon). Aquaculture of Chinook and Coho salmon and Rainbow Trout are major industries in Chile and Australia. Chinook from Chile were released into Argentinean rivers and there were stockings of Coho and Sockeye Salmon and Rainbow Trout in Patagonia.
Species
Some of the species in this genus are highly variable and a number of now-obsolete taxa have been described. In 1989, morphological and genetic studies by Gerald Smith and Ralph Stearley indicated that trouts of the Pacific basin were genetically closer to Pacific salmon (Onchorhynchus species) than to the Salmos – brown trout (Salmo trutta) or Atlantic salmon (Salmo salar) of the Atlantic basin. Thus, in 1989, taxonomic authorities moved the Rainbow, Cutthroat and other Pacific basin trouts into the genus Oncorhynchus.
Extant species
Currently, 12 species and numerous subspecies in this genus are recognized: Behnke (2002).
Fossil species
Based on Stearley & Smith:
†Oncorhynchus belli Stearley & Smith, 2018 (Truckee trout) - late Miocene of Nevada (Truckee Formation). Potentially related to the cutthroat trout.
†Oncorhynchus ketopsis Eiting & Smith, 2007 - Miocene of Oregon (Chalk Hills Formation). Likely belonging to the Pacific salmon group.
†Oncorhynchus lacustris (Cope, 1870) - late Miocene to late Pliocene of Idaho. Likely belonging to the redband trout group.
†Oncorhynchus rastellus Stearley & Smith, 2018 (small spiketooth salmon) - Miocene of Idaho (Chalk Hills Formation). Belongs to the Pacific salmon group.
†Oncorhynchus rastrosus (Cavender & Miller, 1972) (spike-toothed salmon) - middle Miocene to early Pliocene of western North America and Japan. Likely belonging to the Pacific salmon group.
†Oncorhynchus salax Smith, 1975 - Miocene of Idaho (Chalk Hills Formation). Likely belonging to the redband trout group.
| Biology and health sciences | Salmoniformes | Animals |
263745 | https://en.wikipedia.org/wiki/Rack%20and%20pinion | Rack and pinion | A rack and pinion is a type of linear actuator that comprises a circular gear (the pinion) engaging a linear gear (the rack). Together, they convert between rotational motion and linear motion: rotating the pinion causes the rack to be driven in a line. Conversely, moving the rack linearly will cause the pinion to rotate.
The rack and pinion mechanism is used in rack railways, where the pinion mounted on a locomotive or a railroad car engages a rack usually placed between the rails, and helps to move the train up a steep gradient. It is also used in arbor presses and drill presses, where the pinion is connected to a lever and displaces a vertical rack (the ram). In pipelines and other industrial piping systems, a rack displaced by a linar actuator turns a pinion to open or close a valve. Stairlifts, lock gates, electric gates, and the mechanical steering mechanism of cars are other notable applications.
The term "rack and pinion" may be used also when the rack is not straight but arcuate (bent), namely just a section of a large gear.
A single pinion can simultaneously drive two racks, parallel but opposite; which will always be displaced by the same distance, only in opposite directions. Conversely, by applying opposite forces to the two racks one can obtain pure torque on the pinion, without any force component. This double rack and pinion mechanism can be used, for example, with a pair of pneumatic actuators to operate a valve with minimum stress.
Applications gallery
History
The time and place of the invention of the rack-and-pinion mechanism are unknown, but it presumably was not long after the invention of gears. The south-pointing chariot from China and the Antikythera mechanism are evidence of these being well-known already a couple of centuries BC.
In 1598, firearms designer Zhao Shizhen developed the Xuanyuan arquebus (軒轅銃), featuring a rack-and-pinion matchlock mechanism derived from an Ottoman Turkish matchlock design. The Wu Pei Chih (1621) later described Ottoman Turkish muskets that used a rack-and-pinion mechanism.
The use of a variable rack (still using a normal pinion) was invented by Arthur Ernest Bishop in the 1970s, so as to improve vehicle response and steering "feel", especially at high speeds. He also created a low cost press forging process to manufacture the racks, eliminating the need to machine the gear teeth.
Comparison with Worm gear
A rack and pinion has roughly the same purpose as a worm gear with a rack replacing the gear, in that both convert torque to linear force. However the rack and pinion generally provides higher linear speed — since a full turn of the pinion displaces the rack by an amount equal to the pinion's pitch circle whereas a full rotation of the worm screw only displaces the rack by one tooth width. By the same token, a rack and pinion mechanism yields a smaller linear force than a worm gear, for the same input torque. Also, a rack and pinion pair can be used in the opposite way, to turn linear force into torque; whereas a worm drive can be used in only one way.
Geometry
The teeth of a rack and pinion pair may be either straight (parallel to the rotation axis, as in a spur gear) or helical. On the pinion, the profile of the working tooth surfaces is usually an arc of involute, as in most gears. On the rack, on the other hand, the matching working surfaces are flat. One may interpret them as involute tooth faces for a gear with infinite radius. In both parts the teeth are typically formed with a gear cutter (a hob).
| Technology | Mechanisms | null |
263748 | https://en.wikipedia.org/wiki/Scalpel | Scalpel | A scalpel, lancet, or bistoury is a small and extremely sharp bladed instrument used for surgery, anatomical dissection, podiatry and various handicrafts. A lancet is a double-edged scalpel.
Scalpel blades are usually made of hardened and tempered steel, stainless steel, or high carbon steel; in addition, titanium, ceramic, diamond and even obsidian knives are not uncommon. For example, when performing surgery under MRI guidance, steel blades are unusable (the blades would be drawn to the magnets and would also cause image artifacts). Historically, the preferred material for surgical scalpels was silver. Scalpel blades are also offered by some manufacturers with a zirconium nitride–coated edge to improve sharpness and edge retention. Others manufacture blades that are polymer-coated to enhance lubricity during a cut.
Scalpels may be single-use disposable or re-usable. Re-usable scalpels can have permanently attached blades that can be sharpened or, more commonly, removable single-use blades. Disposable scalpels usually have a plastic handle with an extensible blade (like a utility knife) and are used once, then the entire instrument is discarded. Scalpel blades are usually individually packed in sterile pouches but are also offered non-sterile.
Alternatives to scalpels in surgical applications include electrocautery and lasers.
History
Obsidian scalpels older than 2100 BC have been found in a Bronze Age settlement in Turkey. Skulls from the same time and place show signs of brain surgery.
Ancient Egyptians made incisions for embalming with scalpels of sharpened obsidian, a material that is still in use.
The first medical writings of ancient Greeks indicate they were commonly using tools identical to today's scalpels around 500 BC. The amphismela was an anatomical knife-edged on both sides. The term comes from the Greek αμφι (utrinque, "on both sides"), and μελιζω (inside, "I cut").
Ancient Romans used more than 150 different surgical instruments, including scalpels.
10th century Arab-Spanish surgeon Albucasis invented a retractable scalpel.
The French used an amphismela in the 1700s.
South African scientists showed that a blunt scalpel caused sharp cuts if the blade was subjected to ultrasound. Applications might be in energy-saving paper cutting.
Operation
In the palmar grip, also called the "dinner knife" grip, the handle is held with the second through fourth fingers and secured along the base of the thumb, with the index finger extended along the top rear of the blade and the thumb along the side of the handle. This grip is best for initial incisions and larger cuts.
In the pencil grip, best used for more accurate cuts with smaller blades (e.g. #15) and the #7 handle, the scalpel is held with the tips of the first and second fingers and the tip of the thumb with the handle resting on the fleshy base of the index finger and thumb.
Types
Surgical
Surgical scalpels consist of two parts, a blade and a handle. The handles are often reusable, with the blades being replaceable. In medical applications, each blade is only used once (sometimes just for a single, small cut).
The handle is also known as a "B.P. handle", named after Charles Russell Bard and Morgan Parker, founders of the Bard-Parker Company. Morgan Parker patented the 2-piece scalpel design in 1915 and Bard-Parker developed a method of cold sterilization that would not dull the blades, as did the heat-based method that was previously used.
The handle of medical scalpels come in several basic types. The first is a flat handle used in the #3 and #4 handles. The #7 handle is more like a long writing pen, rounded at the front and flat at the back. A #4 handle is larger than a #3. #5 handles are also common, and are round, with a patterning to ensure a non-slip grip. Blades are manufactured with a corresponding fitment size so that they fit on only one size handle. The following table of blades is incomplete and some blades listed may work with handles not specified here.
A lancet has a double-edged blade and a pointed end for making small incisions or drainage punctures.
Handicraft
Graphical and model-making scalpels tend to have round handles, with textured grips (either knurled metal or soft plastic). The blade is usually flat and straight, allowing it to be run easily against a straightedge to produce straight cuts.
There are many kinds of graphic arts blades; the most common around the graphic design studio is the #11 blade which is very similar to a #11 surgical blade (q.v.). Other blade shapes are used for wood carving, cutting leather and heavy fabric.
Blades
Safety
Rising awareness of the dangers of sharps in a medical environment around the beginning of the 21st century led to the development of various methods of protecting healthcare workers from accidental cuts and puncture wounds. According to the Centers for Disease Control and Prevention, as many as 1,000 people were subject to accidental needle sticks and lacerations each day in the United States while providing medical care. Additionally, surgeons can expect to suffer hundreds of such injuries over the course of their career. Scalpel blade injuries were among the most frequent sharps injuries, second only to needlesticks. Scalpel injuries made up 7 percent to 8 percent of all sharps injuries in 2001.
"Scalpel Safety" is a term coined to inform users that there are choices available to them to ensure their protection from this common sharps injury.
Safety scalpels are becoming increasingly popular as their prices come down and also on account of legislation such as the Needle Stick Prevention Act, which requires hospitals to minimize the risk of pathogen transmission through needle or scalpel-related accidents.
There are essentially two kinds of disposable safety scalpels offered by various manufacturers. They can be either classified as retractable blade or retractable sheath type. The retractable blade version made by companies such as OX Med Tech, DeRoyal, Jai Surgicals, Swann Morton, and PenBlade are more intuitive to use due to their similarities to a standard box-cutter. Retractable sheath versions have much stronger ergonomic feel for the doctors and are made by companies such as Aditya Dispomed, Aspen Surgical and Southmedic. A few companies have also started to offer a safety scalpel with a reusable metal handle. In such models, the blade is usually protected in a cartridge. Such systems usually require a custom handle and the price of blades and cartridges is considerably more than for conventional surgical blades.
There are various scalpel blade removers on the market that allows users to safely remove blades from the handle, instead of dangerously using fingers or forceps. In the medical field, when taking into account activation rates, the combination of a single-handed scalpel blade remover with a passing tray or a neutral zone was as safe and up to five times safer than a safety scalpel. There are companies which offer a single-handed scalpel blade remover that complies with regulatory requirements such as US Occupational Safety and Health Administration Standards.
The usage of both safety scalpels and a single-handed blade remover, combined with a hands-free passing technique, are potentially effective in reducing scalpel blade injuries. It is up to employers and scalpel users to consider and use safer and more effective scalpel safety measures when feasible.
| Technology | Surgical instruments | null |
263774 | https://en.wikipedia.org/wiki/Cotinga | Cotinga | The cotingas are a large family, Cotingidae, of suboscine passerine birds found in Central America and tropical South America. Cotingas are birds of forests or forest edges, that are primary frugivorous. They all have broad bills with hooked tips, rounded wings, and strong legs. They range in size from of the fiery-throated fruiteater (Pipreola chlorolepidota) up to of the Amazonian umbrellabird (Cephalopterus ornatus).
Description
Cotingas vary widely in social structure. There is a roughly 50/50 divide in the family between species with biparental care, and those in which the males play no part in raising the young. The purple-throated fruitcrow lives in mixed-sex groups in which one female lays an egg and the others help provide insects to the chick.
In cotinga species where only the females care for the eggs and young, the males have striking courtship displays, often grouped together in leks. Such sexual selection results in the males of these species, including the Guianan cock-of-the-rock, being brightly coloured, or decorated with plumes or wattles, like the umbrellabirds, with their umbrella-like crest and long throat wattles. Other lekking cotingids like the bellbirds and screaming piha, have distinctive and far-carrying calls. In such canopy-dwelling genera as Carpodectes, Cotinga, and Xipholena, males gather high in a single tree or in adjacent trees, but male cocks-of-the-rock, as befits their more terrestrial lives, give their elaborate displays in leks on the ground.
The females of both lekking and biparental species are duller than the males.
Breeding
Nests range from tiny to very large. Many species lay a single egg in a nest so flimsy that the egg can be seen from underneath. This may make the nests hard for predators to find. Fruiteaters build more solid cup nests, and the cocks-of-the-rock attach their mud nests to cliffs. The nests may be open cups or little platforms with loosely woven plant material, usually placed in a tree.
The clutches comprise one to four eggs. Incubation typically takes 15–28 days.
Fledging usually occurs at 28–33 days.
Habitat
Deserts, open woodlands, coastal mangroves, and humid tropical forests comprise their habitats. Cotingas face very serious threats from the loss of their habitats.
Taxonomy and systematics
The family Cotingidae was introduced by French naturalist Charles Lucien Bonaparte in 1849. According to the International Ornithological Committee, as of July 2021, the family contains 66 species divided into 24 genera.
A 2014 molecular phylogenetic study of the cotingas by Jacob Berv and Richard Prum found that the genera formed five monophyletic clades and they proposed that the family could be divided into five subfamilies. The following cladogram is based on a molecular phylogenetic study of the suboscines by Michael Harvey and collaborators published in 2020.
The genus Tijuca was found to be embedded in Lipaugus, a position that was confirmed by a more detailed 2020 study.
A number of species previously placed in this family are now placed in the family Tityridae (genera Laniisoma, Laniocera and Iodopleura)
| Biology and health sciences | Tyranni | Animals |
263800 | https://en.wikipedia.org/wiki/Tadpole | Tadpole | A tadpole or polliwog (also spelled pollywog) is the larval stage in the biological life cycle of an amphibian. Most tadpoles are fully aquatic, though some species of amphibians have tadpoles that are terrestrial. Tadpoles have some fish-like features that may not be found in adult amphibians such as a lateral line, gills and swimming tails. As they undergo metamorphosis, they start to develop functional lungs for breathing air, and the diet of tadpoles changes drastically.
A few amphibians, such as some members of the frog family Brevicipitidae, undergo direct development i.e., they do not undergo a free-living larval stage as tadpoles instead emerging from eggs as fully formed "froglet" miniatures of the adult morphology. Some other species hatch into tadpoles underneath the skin of the female adult or are kept in a pouch until after metamorphosis. Having no hard skeletons, it might be expected that tadpole fossils would not exist. However, traces of biofilms have been preserved and fossil tadpoles have been found dating back to the Middle Jurassic.
Tadpoles are eaten as human food in some parts of the world and are mentioned in various folk tales from around the world.
Etymology
The name tadpole is from Middle English , made up of the elements , 'toad', and , 'head' (modern English poll). Similarly, pollywog / polliwog is from Middle English , made up of the same , 'head', and , 'to wiggle'.
General description
The life cycle of all amphibians involves a larval stage that is intermediate between embryo and adult. In most cases this larval stage is a limbless free-living organism that has a tail and is referred to as a tadpole, although in a few cases (e.g., in the Breviceps and Probreviceps genera of frogs) direct development occurs in which the larval stage is confined within the egg. Tadpoles of frogs are mostly herbivorous, while tadpoles of salamanders and caecilians are carnivorous.
Anura
Tadpoles of frogs and toads are usually globular, with a laterally compressed tail with which they swim by lateral undulation. When first hatched, anuran tadpoles have external gills that are eventually covered by skin, forming an opercular chamber with internal gills vented by spiracles. Depending on the species, there can be two spiracles on both sides of the body, a single spiracle on the underside near the vent, or a single spiracle on the left side of the body. Newly hatched tadpoles are also equipped with a cement gland which allows them to attach to objects. The tadpoles have a cartilaginous skeleton and a notochord which eventually develops into a proper spinal cord.
Anuran tadpoles are usually herbivorous, feeding on soft decaying plant matter. The gut of most tadpoles is long and spiral-shaped to efficiently digest organic matter and can be seen through the bellies of many species. Though many tadpoles will feed on dead animals if available to them, only a few species of frog have strictly carnivorous tadpoles, an example being the frogs of the family Ceratophryidae, their cannibalistic tadpoles having wide gaping mouths with which they devour other organisms, including other tadpoles. Another example is the tadpoles of the New Mexico spadefoot toad (Spea multiplicata) which will develop a carnivorous diet along with a broader head, larger jaw muscles, and a shorter gut if food is scarce, allowing them to consume fairy shrimp and their smaller herbivorous siblings. A few genera such as Pipidae and Microhylidae have species whose tadpoles are filter feeders that swim through the water column feeding on plankton. Megophrys tadpoles feed at the water surface using unusual funnel-shaped mouths.
As a frog tadpole matures it gradually develops its limbs, with the back legs growing first and the front legs second. The tail is absorbed into the body using apoptosis. Lungs develop around the time as the legs start growing, and tadpoles at this stage will often swim to the surface and gulp air. During the final stages of metamorphosis, the tadpole's mouth changes from a small, enclosed mouth at the front of the head to a large mouth the same width as the head. The intestines shorten as they transition from a herbivorous diet to the carnivorous diet of adult frogs.
Tadpoles vary greatly in size, both during their development and between species. For example, in a single family, Megophryidae, length of late-stage tadpoles varies between and . The tadpoles of the paradoxical frog (Pseudis paradoxa) can reach up to , the longest of any frog, before shrinking to a mere snout-to-vent length of 3.4–7.6 cm (1.3–3.0 in).
While most anuran tadpoles inhabit wetlands, ponds, vernal pools, and other small bodies of water with slow moving water, a few species are adapted to different environments. Some frogs have terrestrial tadpoles, such as the family Ranixalidae, whose tadpoles are found in wet crevices near streams. The tadpoles of Micrixalus herrei are adapted to a fossorial lifestyle, with a muscular body and tail, eyes covered by a layer of skin, and reduced pigment. Several frogs have stream dwelling tadpoles equipped with a strong oral sucker that allows them to hold onto rocks in fast flowing water, two examples being the Indian purple frog (Nasikabatrachus sahyadrensis) and the tailed frogs (Ascaphus) of Western North America. Although there are no marine tadpoles, the tadpoles of the crab-eating frog can cope with brackish water.
Some anurans will provide parental care towards their tadpoles. Frogs of the genus Afrixalus will lay their eggs on leaves above water, folding the leaves around the eggs for protection. Female Pipa frogs will embed the eggs into their backs where they get covered by a thin layer of skin. The eggs will hatch underneath her skin and grow, eventually leaving as either large tadpoles (such as in Pipa parva) or as fully formed froglets (Pipa pipa). Female marsupial frogs (Hemiphractidae) will carry eggs on her back for various amounts of time, with it going as far as letting the tadpoles develop into tiny froglets in a pouch. Male African bullfrogs (Pyxicephalus adspersus) will keep watch over their tadpoles, attacking anything that might be a potential threat, even though he may eat some of the tadpoles himself.
Males of the Emei mustache toads (Leptobrachium boringii) will construct nests along riverbanks where they breed with females and keep watch over the eggs, losing as much as 7.3% of their body mass in the time they spend protecting the nest. Male midwife toads (Alytes) will carry eggs between their legs to protect them from predators, eventually releasing them into a body of water when they are ready to hatch. Poison dart frogs (Dendrobatidae) will carry their tadpoles to various locations, usually phytotelma, where they remain until metamorphosis. Some female dart frogs such as the strawberry poison dart frog (Oophaga pumilio) will regularly lay unfertilized eggs for the developing tadpoles to feed on.
Fossil record
Despite their soft-bodied nature and lack of mineralised hard parts, fossil tadpoles (around 10 cm in length) have been recovered from Upper Miocene strata. They are preserved by virtue of biofilms, with more robust structures (the jaw and bones) preserved as a carbon film. In Miocene fossils from Libros, Spain, the brain case is preserved in calcium carbonate, and the nerve cord in calcium phosphate. Other parts of the tadpoles' bodies exist as organic remains and bacterial biofilms, with sedimentary detritus present in the gut. Tadpole remains with telltale external gills are also known from several labyrinthodont groups. The oldest unambiguous fossil, a tadpole of the species Notobatrachus deguistioi from the Middle Jurassic, was published on in 2024.
Human use
Tadpoles are used in a variety of cuisines. Tadpoles of the megophryid frog Oreolalax rhodostigmatus are particularly large, more than in length, and are collected for human consumption in China. In Peru Telmatobius mayoloi tadpoles are collected for both food and medicine.
Mythology and history
According to Sir George Scott, in the origin myths of the Wa people in China and Myanmar, the first Wa originated from two female ancestors Ya Htawm and Ya Htai, who spent their early phase as tadpoles ("") in a lake in the Wa country known as Nawng Hkaeo.
| Biology and health sciences | Animal ontogeny | null |
263886 | https://en.wikipedia.org/wiki/Differential%20operator | Differential operator | In mathematics, a differential operator is an operator defined as a function of the differentiation operator. It is helpful, as a matter of notation first, to consider differentiation as an abstract operation that accepts a function and returns another function (in the style of a higher-order function in computer science).
This article considers mainly linear differential operators, which are the most common type. However, non-linear differential operators also exist, such as the Schwarzian derivative.
Definition
Given a nonnegative integer m, an order- linear differential operator is a map from a function space on to another function space that can be written as:
where is a multi-index of non-negative integers, , and for each , is a function on some open domain in n-dimensional space. The operator is interpreted as
Thus for a function :
The notation is justified (i.e., independent of order of differentiation) because of the symmetry of second derivatives.
The polynomial p obtained by replacing partials by variables in P is called the total symbol of P; i.e., the total symbol of P above is:
where The highest homogeneous component of the symbol, namely,
is called the principal symbol of P. While the total symbol is not intrinsically defined, the principal symbol is intrinsically defined (i.e., it is a function on the cotangent bundle).
More generally, let E and F be vector bundles over a manifold X. Then the linear operator
is a differential operator of order if, in local coordinates on X, we have
where, for each multi-index α, is a bundle map, symmetric on the indices α.
The kth order coefficients of P transform as a symmetric tensor
whose domain is the tensor product of the kth symmetric power of the cotangent bundle of X with E, and whose codomain is F. This symmetric tensor is known as the principal symbol (or just the symbol) of P.
The coordinate system xi permits a local trivialization of the cotangent bundle by the coordinate differentials dxi, which determine fiber coordinates ξi. In terms of a basis of frames eμ, fν of E and F, respectively, the differential operator P decomposes into components
on each section u of E. Here Pνμ is the scalar differential operator defined by
With this trivialization, the principal symbol can now be written
In the cotangent space over a fixed point x of X, the symbol defines a homogeneous polynomial of degree k in with values in .
Fourier interpretation
A differential operator P and its symbol appear naturally in connection with the Fourier transform as follows. Let ƒ be a Schwartz function. Then by the inverse Fourier transform,
This exhibits P as a Fourier multiplier. A more general class of functions p(x,ξ) which satisfy at most polynomial growth conditions in ξ under which this integral is well-behaved comprises the pseudo-differential operators.
Examples
The differential operator is elliptic if its symbol is invertible; that is for each nonzero the bundle map is invertible. On a compact manifold, it follows from the elliptic theory that P is a Fredholm operator: it has finite-dimensional kernel and cokernel.
In the study of hyperbolic and parabolic partial differential equations, zeros of the principal symbol correspond to the characteristics of the partial differential equation.
In applications to the physical sciences, operators such as the Laplace operator play a major role in setting up and solving partial differential equations.
In differential topology, the exterior derivative and Lie derivative operators have intrinsic meaning.
In abstract algebra, the concept of a derivation allows for generalizations of differential operators, which do not require the use of calculus. Frequently such generalizations are employed in algebraic geometry and commutative algebra. | Mathematics | Differential calculus | null |
263902 | https://en.wikipedia.org/wiki/Evanescent%20field | Evanescent field | In electromagnetics, an evanescent field, or evanescent wave, is an oscillating electric and/or magnetic field that does not propagate as an electromagnetic wave but whose energy is spatially concentrated in the vicinity of the source (oscillating charges and currents). Even when there is a propagating electromagnetic wave produced (e.g., by a transmitting antenna), one can still identify as an evanescent field the component of the electric or magnetic field that cannot be attributed to the propagating wave observed at a distance of many wavelengths (such as the far field of a transmitting antenna).
A hallmark of an evanescent field is that there is no net energy flow in that region. Since the net flow of electromagnetic energy is given by the average Poynting vector, this means that the Poynting vector in these regions, as averaged over a complete oscillation cycle, is zero.
Use of the term
In many cases one cannot simply say that a field is or is not "evanescent" – having the Poynting vector average to zero in some direction (or all directions).
In most cases where they exist, evanescent fields are simply thought of and referred to the same as all other electric or magnetic fields involved, without any special recognition of those fields' evanescence. The term's use is mostly limited to distinguishing a part of a field or solution in those cases where one might only expect the fields of a propagating wave.
For instance, in the illustration at the top of the article, energy is indeed carried in the horizontal direction. However, in the vertical direction, the field strength drops off exponentially with increasing distance above the surface. This leaves most of the field concentrated in a thin boundary layer very close to the interface; for that reason, it is referred to as a surface wave. However, despite energy flowing horizontally, along the vertical there is no net propagation of energy away from (or toward) the surface, so that one could properly describe the field as being "evanescent in the vertical direction". This is one example of the context dependence of the term.
Everyday electronic devices and electrical appliances are surrounded by large fields which are evanescent; their operation involves alternating voltages (producing an electric field between them) and alternating currents (producing a magnetic field around them) which are expected to only carry power along internal wires, but not to the outsides of the devices. Even though the term "evanescent" is not mentioned in this ordinary context, the appliances' designers still may be concerned with maintaining evanescence, in order to prevent or limit production of a propagating electromagnetic wave, which would lead to radiation loss, since a propagating wave "steals" its power from the circuitry or donates unwanted interference.
The term "evanescent field" does arise in various contexts where a propagating electromagnetic wave is involved (even if confined). The term then differentiates electromagnetic field components that accompany the propagating wave, but which do not themselves propagate. In other, similar cases, where a propagating electromagnetic wave would normally be expected (such as light refracted at the interface between glass and air), the term is invoked to describe that part of the field where the wave is suppressed (such as light traveling through glass, impinging on a glass-to-air interface but beyond the critical angle).
Although all electromagnetic fields are classically governed according to Maxwell's equations, different technologies or problems have certain types of expected solutions, and when the primary solutions involve wave propagation the term evanescent is frequently applied to field components or solutions which do not share that property.
For instance, the propagation constant of a hollow metal waveguide is a strong function of frequency (a dispersion relation). Below a certain frequency (the cut-off frequency) the propagation constant becomes an imaginary number. A solution to the wave equation having an imaginary wavenumber does not propagate as a wave but falls off exponentially, so the field excited at that lower frequency is considered evanescent. It can also be simply said that propagation is "disallowed" for that frequency.
The formal solution to the wave equation can describe modes having an identical form, but the change of the propagation constant from real to imaginary as the frequency drops below the cut-off frequency totally changes the physical nature of the result. The solution may be described as a "cut-off mode" or an "evanescent mode"; while a different author will just state that no such mode exists. Since the evanescent field corresponding to the mode was computed as a solution to the wave equation, it is often discussed as being an "evanescent wave" even though its properties (such as carrying no energy) are inconsistent with the definition of wave.
Although this article concentrates on electromagnetics, the term evanescent is used similarly in fields such as acoustics and quantum mechanics, where the wave equation arises from the physics involved. In these cases, solutions to the wave equation resulting in imaginary propagation constants are likewise called "evanescent", and have the essential property that no net energy is transferred, even though there is a non-zero field.
Evanescent wave applications
In optics and acoustics, evanescent waves are formed when waves traveling in a medium undergo total internal reflection at its boundary because they strike it at an angle greater than the critical angle. The physical explanation for the existence of the evanescent wave is that the electric and magnetic fields (or pressure gradients, in the case of acoustical waves) cannot be discontinuous at a boundary, as would be the case if there was no evanescent wave field. In quantum mechanics, the physical explanation is exactly analogous—the Schrödinger wave-function representing particle motion normal to the boundary cannot be discontinuous at the boundary.
Electromagnetic evanescent waves have been used to exert optical radiation pressure on small particles to trap them for experimentation, or to cool them to very low temperatures, and to illuminate very small objects such as biological cells or single protein and DNA molecules for microscopy (as in the total internal reflection fluorescence microscope). The evanescent wave from an optical fiber can be used in a gas sensor, and evanescent waves figure in the infrared spectroscopy technique known as attenuated total reflectance.
In electrical engineering, evanescent waves are found in the near-field region within one third of a wavelength of any radio antenna. During normal operation, an antenna emits electromagnetic fields into the surrounding nearfield region, and a portion of the field energy is reabsorbed, while the remainder is radiated as EM waves.
Recently, a graphene-based Bragg grating (one-dimensional photonic crystal) has been fabricated and demonstrated its competence for excitation of surface electromagnetic waves in the periodic structure using a prism coupling technique.
In quantum mechanics, the evanescent-wave solutions of the Schrödinger equation give rise to the phenomenon of wave-mechanical tunneling.
In microscopy, systems that capture the information contained in evanescent waves can be used to create super-resolution images. Matter radiates both propagating and evanescent electromagnetic waves. Conventional optical systems capture only the information in the propagating waves and hence are subject to the diffraction limit. Systems that capture the information contained in evanescent waves, such as the superlens and near field scanning optical microscopy, can overcome the diffraction limit; however these systems are then limited by the system's ability to accurately capture the evanescent waves. The limitation on their resolution is given by
where is the maximal wave vector that can be resolved, is the distance between the object and the sensor, and is a measure of the quality of the sensor.
More generally, practical applications of evanescent waves can be classified as (1) those in which the energy associated with the wave is used to excite some other phenomenon within the region of space where the original traveling wave becomes evanescent (for example, as in the total internal reflection fluorescence microscope) or (2) those in which the evanescent wave couples two media in which traveling waves are allowed, and hence permits the transfer of energy or a particle between the media (depending on the wave equation in use), even though no traveling-wave solutions are allowed in the region of space between the two media. An example of this is wave-mechanical tunnelling, and is known generally as evanescent wave coupling.
Total internal reflection of light
For example, consider total internal reflection in two dimensions, with the interface between the media lying on the x axis, the normal along y, and the polarization along z. One might expect that for angles leading to total internal reflection, the solution would consist of an incident wave and a reflected wave, with no transmitted wave at all, but there is no such solution that obeys Maxwell's equations. Maxwell's equations in a dielectric medium impose a boundary condition of continuity for the components of the fields E||, H||, Dy, and By. For the polarization considered in this example, the conditions on E|| and By are satisfied if the reflected wave has the same amplitude as the incident one, because these components of the incident and reflected waves superimpose destructively. Their Hx components, however, superimpose constructively, so there can be no solution without a non-vanishing transmitted wave. The transmitted wave cannot, however, be a sinusoidal wave, since it would then transport energy away from the boundary, but since the incident and reflected waves have equal energy, this would violate conservation of energy. We therefore conclude that the transmitted wave must be a non-vanishing solution to Maxwell's equations that is not a traveling wave, and the only such solutions in a dielectric are those that decay exponentially: evanescent waves.
Mathematically, evanescent waves can be characterized by a wave vector where one or more of the vector's components has an imaginary value. Because the vector has imaginary components, it may have a magnitude that is less than its real components.
For the plane of incidence as the plane at and the interface of the two mediums as the plane at , the wave vector of the transmitted wave has the form
with and , where is the magnitude of the wave vector of the transmitted wave (so the wavenumber), is the angle of refraction, and and are the unit vectors along the axis direction and the axis direction respectively.
By using the Snell's law where , , and are the refractive index of the medium where the incident wave and the reflected wave exist, the refractive index of the medium where the transmitted wave exists, and the angle of incidence respectively,
.
with .
If a part of the condition of the total internal reflection as , is satisfied, then
.
If the polarization is perpendicular to the plane of incidence (along the direction), then the electric field of any of the waves (incident, reflected, or transmitted) can be expressed as
where is the unit vector in the axis direction.
By assuming plane waves as , and substituting the transmitted wave vector into , we find for the transmitted wave:
where is the attenuation constant, and is the phase constant. is ignored since it does not physically make sense (the wave amplification along y the direction in this case).
Evanescent-wave coupling
Especially in optics, evanescent-wave coupling refers to the coupling between two waves due to physical overlap of what would otherwise be described as the evanescent fields corresponding to the propagating waves.
One classical example is frustrated total internal reflection (FTIR) in which the evanescent field very close (see graph) to the surface of a dense medium at which a wave normally undergoes total internal reflection overlaps another dense medium in the vicinity. This disrupts the totality of the reflection, diverting some power into the second medium.
Coupling between two optical waveguides may be effected by placing the fiber cores close together so that the evanescent field generated by one element excites a wave in the other fiber. This is used to produce fiber-optic splitters and in fiber tapping. At radio (and even optical) frequencies, such a device is called a directional coupler. The device is usually called a power divider in the case of microwave transmission and modulation.
Evanescent-wave coupling is synonymous with near field interaction in electromagnetic field theory. Depending on the nature of the source element, the evanescent field involved is either predominantly electric (capacitive) or magnetic (inductive), unlike (propagating) waves in the far field where these components are connected (identical phase, in the ratio of the impedance of free space). The evanescent wave coupling takes place in the non-radiative field near each medium and as such is always associated with matter; i.e., with the induced currents and charges within a partially reflecting surface. In quantum mechanics the wave function interaction may be discussed in terms of particles and described as quantum tunneling.
Applications
Evanescent wave coupling is commonly used in photonic and nanophotonic devices as waveguide sensors or couplers (see e.g., prism coupler).
Evanescent wave coupling is used to excite, for example, dielectric microsphere resonators.
Evanescent coupling, as near field interaction, is one of the concerns in electromagnetic compatibility.
Coupling of optical fibers without loss for fiber tapping.
Evanescent wave coupling plays a major role in the theoretical explanation of extraordinary optical transmission.
Evanescent wave coupling is used in powering devices wirelessly.
A total internal reflection fluorescence microscope uses the evanescent wave produced by total internal reflection to excite fluorophores close to a surface. This is useful when surface properties of biological samples need to be studied.
| Physical sciences | Electromagnetic radiation | Physics |
263960 | https://en.wikipedia.org/wiki/Warship | Warship | A warship or combatant ship is a ship that is used for naval warfare. Usually they belong to the navy branch of the armed forces of a nation, though they have also been operated by individuals, cooperatives and corporations. As well as being armed, warships are designed to withstand damage and are typically faster and more maneuverable than merchant ships. Unlike a merchant ship, which carries cargo, a warship typically carries only weapons, ammunition and supplies for its crew.
In wartime, the distinction between warships and merchant ships is often blurred. Until the 17th century it was common for merchant ships to be pressed into naval service, and not unusual for more than half of a fleet to be composed of merchant ships—there was not a large difference in construction, unlike the difference between a heavily armoured battleship and an ocean liner. Until the threat of piracy subsided in the 19th century, it was normal practice to arm larger merchant ships such as galleons. Warships have also often been used as troop carriers or supply ships, such as by the French Navy in the 18th century or the Imperial Japanese Navy during the Second World War. In war since the early 20th century, merchant ships were often armed and used as auxiliary warships, such as the Q-ships of the First World War and the armed merchant cruisers of the Second World War.
Types and classes
The main types of warships today are, in order of decreasing size: aircraft carriers – amphibious assault ships – cruisers – destroyers – frigates – corvettes – fast attack boats. A more extensive list follows:
Submarine, are ships capable of staying submerged for days. Modern submarines can stay underwater for months, with food supplies as the only limiting factor.
Fleet submarine is a type of submarine with the speed, range, and endurance to operate as part of a navy's battle fleet.
Ballistic Missile Submarine is a submarine capable of deploying submarine-launched ballistic missiles (SLBMs) with nuclear warheads.
Cruise missile submarine are submarines equipped with cruise missiles.
Attack Submarine is a submarine with the purpose of attacking other submarines.
Coastal submarine or littoral submarine is a small, maneuverable type of submarine with shallow draft well suited to navigation of coastal channels and harbors.
Midget submarine is any submarine under 150 tons, typically operated by a crew of one or two but sometimes up to six or nine.
Submarine aircraft carrier is a submarine equipped with aircraft for observation or attack missions. These submarines saw their most extensive use during World War II, although their operational significance remained small.
Cruiser submarine were a type of a very large submarine designed to remain at sea for extended periods in areas distant from base facilities.
Amphibious warfare ships are warships employed to land and support ground forces, such as marines, on enemy territory during an amphibious assault.
Amphibious assault ship is a type of amphibious warfare ship employed to land and support ground forces with ship-deployed helicopters and V/STOL aircraft on enemy territory by an amphibious assault.
Landing helicopter dock is a multipurpose amphibious assault ship that is capable of operating as a helicopter carrier and also has a well deck for supporting landing crafts.
Amphibious transport dock is an amphibious warfare ship, that embarks, transports, and lands elements of a landing force for expeditionary warfare missions.
Dock landing ship is an amphibious warfare ship with a well dock to transport and launch landing craft and amphibious vehicles.
Landing craft are small and medium seagoing watercraft, such as boats and barges, used to convey a landing force (infantry and vehicles) from the sea to the shore during an amphibious assault.
Landing Craft Utility is a type of landing craft used by amphibious forces to transport equipment and troops to the shore. They are capable of transporting tracked or wheeled vehicles and marines from amphibious assault ships to beachheads.
Landing Craft Mechanized is a landing craft designed for carrying vehicles during amphibious assaults.
Landing ship, tank is the naval designation for ships first developed during World War II to support amphibious operations by carrying tanks, armoured fighting vehicles, transport vehicles, cargo, and landing troops directly onto shore with no docks or piers.
Landing Craft Support were two distinct classes of amphibious warfare vessels used by the United States Navy during World War II to support landing crafts.
Capital ship, the largest and most important ships in a nation's fleet. These were previously battlecruisers, battleships, and aircraft carriers, but the first two warship types are now no longer used.
Aircraft carrier, a warship primarily armed with carrier-based aircraft.
Fleet carrier is an aircraft carrier designed to operate with the main fleet of a nation's navy.
Light aircraft carrier is an aircraft carrier that is smaller than the standard carriers of a navy.
Escort carrier, also called a "jeep carrier" or "baby flattop" is a slow type of aircraft carrier used during WWII.
Anti-submarine warfare carrier is a type of small aircraft carrier whose primary role is as the nucleus of an anti-submarine warfare hunter-killer group.
Aircraft cruiser (also known as aviation cruiser or cruiser-carrier) is a type of warship that combines the features of the aircraft carrier and a surface warship, such as a cruiser or battleship.
Helicopter cruiser.
Helicopter carrier, an aircraft carrier especially suited to carry helicopters and V/STOL aircraft.
Seaplane tender, a boat or ship that supports the operation of seaplanes. Some of these vessels, also known as seaplane carriers, could not only carry seaplanes but also provided all the facilities needed for their operation; these ships are regarded by some as the first aircraft carriers and became obsolete at the end of the Second World War.
Battleship, a large, heavily armored warship equipped with many powerful guns. A term which generally post-dates sailing warships.
Ironclad battleship, battleships built before the pre-dreadnought in the late 1850s to the early 1890s.
Pre-dreadnought battleship, sea-going battleships built to a common design before the launch of dreadnoughts, between the mid-1880s to the early 1900s. Pre-dreadnoughts commonly featured a mixed main battery composed of several different caliber guns.
Dreadnought, an early 20th-century battleship, which set the pattern for all subsequent battleship construction. Dreadnoughts differ from pre-dreadnoughts in that they feature an all-big-gun main battery. The advantage lies in that if all the big guns have the same characteristics, only one firing solution will be needed to aim them all.
Fast battleships were battleships which emphasized speed without – in concept – undue compromise of either armor or armament.
Battlecruiser, a ship with battleship-level armament and cruiser-level armour; typically faster than a battleship because the reduction in armour allowed mounting of more powerful propulsion machinery, or the use of a more slender hull shape with a lower drag coefficient.
Cruiser, a fast, independent warship. Traditionally, cruisers were the smallest warships capable of independent action. As of 2024, only two countries operated active duty vessels formally classed as cruisers: the United States and Russia.
Guided missile cruisers are cruisers armed anti-ship missiles.
Torpedo cruiser is a type of cruiser that is armed primarily with torpedoes.
Armored cruiser was a type of warship of the late 19th and early 20th centuries. It was designed like other types of cruisers to operate as a long-range, independent warship, capable of defeating any ship apart from a battleship and fast enough to outrun any battleship it encountered.
Large cruiser is the class of the battlecruiser-sized Alaska-class cruisers of the United States Navy during World War II.
Heavy cruiser was a type of cruiser, a naval warship designed for long range and high speed, armed generally with naval guns of roughly 203 mm (8 in) caliber; its parameters were dictated by the Washington Naval Treaty of 1922 and the London Naval Treaty of 1930.
Pocket battleship, nickname for the Deutschland-class heavy cruisers.
Light cruiser is a type of small or medium-sized warship. The term is a shortening of the phrase "light armored cruiser", describing a small ship that carried armor in the same way as an armored cruiser: a protective belt and deck.
Scout cruiser was a type of warship of the early 20th century, which were smaller, faster, more lightly armed and armoured than protected cruisers or light cruisers, but larger than contemporary destroyers. They were used for scouting.
Protected cruiser is a type of naval cruiser of the late-19th century, gained their description because an armored deck offered protection for vital machine-spaces from fragments caused by shells exploding above them.
Unprotected cruiser was a type of naval cruiser in use during the early 1870s Victorian or pre-dreadnought era.
Coastal defence ship were a type of cruiser-sized warship built for the purpose of coastal defense.
Destroyer, a fast and highly maneuverable warship, traditionally incapable of independent action. Originally developed to counter the threat of torpedo boats, they are now the largest independent warship generally seen on the ocean.
Guided missile destroyer are destroyers armed with anti-ship missiles.
Escort destroyer was a small warship built to full naval standards which was optimised for air-defence and anti-submarine duties in wartime, but which retained many of the capabilities of a traditional fleet destroyer, enabling it to conduct operations in conjunction with main fleet units as well as carrying out convoy escort and ASW patrols.
Destroyer escort was the United States Navy mid-20th-century classification for a 20-knot (37 km/h; 23 mph), warship designed with the endurance necessary to escort mid-ocean convoys of merchant marine ships similar to frigates.
Frigate, originally a medium-sized sailing ship. Although they date back to the 17th century, frigates in modern navies are typically used to protect merchant ships and other warships.
Armoured frigate are frigates with armour which was added to ships based on existing frigate and ship of the line designs. The additional weight of the armour on these first ironclad warships meant that they could have only one gun deck, and they were technically frigates, even though they were more powerful than existing ships-of-the-line and occupied the same strategic role.
Guided missile frigates are frigates armed with anti-ship missiles.
Corvettes were small ships during the age of sail. The concept was revived again in WWII as a merchant convoy escort and anti-submarine ship. Today they are used for anti-submarine warfare and patrolling.
Littoral Combat Ship is a United States Navy classification of warships with the size and role of corvettes.
Fast attack crafts are a small, fast, agile, offensive, often affordable type of warships armed with anti-ship missiles, guns or torpedoes.
Missile boats are small, fast warship armed with anti-ship missiles.
Torpedo boat, are small, fast surface vessels designed for launching torpedoes.
Torpedo ram is a type of torpedo boat combining a ram with torpedo tubes.
Motor torpedo boat is a type of fast torpedo boat, especially of the mid 20th century.
Patrol vessels are relatively small naval vessels generally designed for coastal defence, border protection, immigration law-enforcement, search and rescue duties. They may be broadly classified as inshore patrol vessels or offshore patrol vessels.
Mine warfare vessels:
Minesweeper are small warships designed to remove or detonate naval mines.
Minehunter are naval vessels that seek, detect, and destroys individual naval mines.
Mine countermeasures vessels are atype of naval ships designed for the location of and destruction of naval mines which combines the role of a minesweeper and minehunter in one hull.
Minelayer are naval vessels that plant naval mines offshore.
Fire ship, a vessel of any sort set on fire and sent into an anchorage or fleet with the intention of causing destruction and chaos. Exploding fire ships are called hellburners. The development of unmanned surface vehicles has revived the use of fire ships.
Naval drifters are boats built along the lines of a commercial fishing drifter but fitted out for naval purposes.
Naval trawlers are vessels built along the lines of fishing trawlers but fitted out for naval purposes.
Armed merchantman is a type of merchant ship equipped with naval guns, usually for defensive purposes, either by design or after the fact.
Commerce raider, any armed vessel—privately or government-owned—sanctioned to raid a nation's merchant fleet.
Merchant raiders are disguised commerce raiders.
Gunboats are naval watercraft designed for the express purpose of carrying one or more guns to bombard coastal targets.
River gunboat is a type of gunboat for riverine use.
Flat-iron gunboats were a number of classes of coastal gunboats generally characterized by small size, low freeboard, the absence of masts, and the mounting of a single non-traversing large gun, aimed by pointing the vessel.
Torpedo gunboat were a form of gunboat armed with torpedoes and designed for hunting and destroying smaller torpedo boats.
Motor gunboat is a type of a fast gunboat armed with machine-guns and autocannons.
Monitor, a type of small, heavily gunned warships with shallow draft designed for shore bombardment.
River monitor, a type of monitors used in rivers.
Breastwork monitor was a modification of the monitor by Sir Edward Reed of the Royal Navy.
Q-ship, also known as Q-boats, decoy vessels, special service ships, or mystery ships, were heavily armed merchant ships with concealed weaponry, designed to lure submarines into making surface attacks. This gave Q-ships the chance to open fire and sink them.
Submarine chaser was a small warship used in anti-submarine warfare.
Armed yachts were modified yachts that were armed with weapons and were typically in the service of a navy.
Balloon carrier was a type of ship equipped with a hot-air balloon tied to the ship with a rope or cable, which was used for observation. This type of ship was later replaced by seaplane tenders and aircraft carriers.
Sloop-of-war was a sailing vessel category later revived in WWII as a convoy escort ship.
Screw sloop was a propeller-driven sloop-of-wars used during the mid-19th Century.
Ironclad, a wooden warship with external iron plating.
Casemate ironclad were a type of ironclad gunboats used in the American Civil War.
Central battery ship in European continental navies, was a development of the (high-freeboard) broadside ironclad of the 1870s
Turret ship was a 19th-century type of warship, the earliest to have their guns mounted in a revolving gun turret, instead of a broadside arrangement.
Floating battery is a kind of armed watercraft, often improvised or experimental, which carries heavy armament but has few other qualities as a warship.
Ship of the line, a sailing warship capable of standing in the line of battle. A direct predecessor to the later battleship.
Cottonclad warships were steam-powered warships with bales of cotton lining as armour used in the American Civil War. The armaments consisted of a ram, random numbers of different cannons and sharpshooters.
Brig of War is a brig armed for use by a navy.
Bomb vessels were wooden sailing ships which carried mortars instead of cannons.
Dispatch boats were small boats, and sometimes large ships, tasked to carry military dispatches from ship to ship or from ship to shore or, in some cases from shore to shore. Dispatch boats were employed when other means of transmitting a message was not possible or safe or as quick.
Aviso, a kind of dispatch boat.
Man-of-war, a British Navy expression for a sailing warship.
Grab was a type of ship common on the Malabar Coast in the 18th and 19th centuries. The ghurāb was originally a galley, but the type evolved into sailing ships armed with cannons.
Gallivat were small, armed type of boats, with sails and oars, armed with swivel guns and used on the Malabar Coast in the 18th and 19th centuries.
Galleass, a sailing and rowing warship, equally well suited to sailing and rowing.
Galleon, a 16th-century armed cargo carriers.
Galley was a type of ship optimised for propulsion by oars. Galleys were historically used for warfare, trade, and piracy.
War canoe was a kind of watercraft of the canoe type designed and outfitted for warfare using bow, spear and shield wielding warriors. During the gunpowder era a single brass or iron cannon was mounted on the bow or stern along with musketeers. These warships were used by many tribes and cultures all around the globe.
Longship, a Viking raiding ship.
East Asian warships:
Geobukseon (literally Turtle ship) were wooden sail and oar propelled Korean warships armed with cannons.
Panokseon (literally board roofed ship) were a type of Korean wooden warships propelled by both sailing and rowing armed with cannons and Hwacha multiple rocket launchers.
Atakebune were wooden oar propelled 16th Century Japanese warships armed with few cannons, arquebusiers, and archers. They were mostly bulky floating fortifications.
Mengchong (literally Covered Assaulter) was a type of leather-covered assault warship used in the 2nd and 3rd centuries CE in China.
Louchuan (literally Tower Ship) was a type of warship used as a floating fortress in Ancient China. The Louchuan was meant to board troops onto enemy ships. Although they were also armed with trebuchets for ranged combat.
Wugongchuan (literally Centipede Ship) was a Chinese oared vessel of the 16th century inspired by the Portuguese galley.
Hellenistic galleys, warships propelled by oars with a sail for use in favorable winds used in the Mediterranean Sea:
Bireme, an ancient vessel, propelled by two banks of oars.
Trireme, an ancient warship propelled by three banks of oars.
Quadrireme, an ancient warship invented in Carthage with two levels of oarsmen, and was therefore lower than the quinquereme.
Quinquereme, an ancient warship propelled by three banks of oars. On the upper row, two rowers hold one oar; on the middle row, two rowers; and on the lower row, one man to an oar.
Hexareme, an ancient warship invented in Syracuse. The exact arrangement of the hexareme's oars is unclear. If it evolved naturally from the earlier designs, it would be a trireme with two rowers per oar.
Septireme, an ancient warship invented by the Macedonia, the septireme was derived by adding a standing rower to the lower level of the hexareme.
Octeres, very little is known about the octeres, at least two of their type were in the fleet of Philip V of Macedon at the Battle of Chios.
Enneres, a type of warship whose oaring system may have been a modification of the quadrireme, with two teams of five and four oarsmen.
Deceres, a type of warship which is present alongside "nines" in the fleet of Antigonus I Monophthalmus in 315 BC. It is most likely that the "ten" was derived from adding another oarsman to the "nine".
Tessarakonteres, a very large catamaran galley reportedly built by Ptolemy IV Philopator of Egypt. It had seven naval rams, with one primary, and the deck would have provided a stable platform for catapults that were often mounted on supergalleys. However, the "forty" was likely just a showpiece; Plutarch describes the ship as for exhibition only.
Lembos, light warships most commonly associated with the vessels used by the Illyrian tribes, chiefly for piracy, in the area of Dalmatia. Was soon adopted by Macedonia, Seleucid Empire, Roman Republic and Sparta.
Hemiolia, light and fast warship that appeared in the early 4th century BC. It was particularly favoured by pirates in the eastern Mediterranean, but also used by Alexander the Great as far as the rivers Indus and Jhelum, and by the Romans as a troop transport. According to one view, it was manned by half the number of oarsmen to make room for the soldiers. According to another, there were one and a half files of oarsmen on each side, with the additional half file placed amidships, where the hull was wide enough to accommodate them.
Trihemiolia, this type was classed with the trireme, and had two and a half files of oarsmen on each side. Judging from the Lindos relief and the famous Nike of Samothrace, both of which are thought to represent trihemioliai, the two upper files would have been accommodated in an oarbox, with the half-file located beneath them in the classic thalamitai position of the trireme.
Liburna, a type of small galley used for raiding and patrols. It was originally used by the Liburnians, a pirate tribe from Dalmatia, and later used by the Roman Navy. It had one bench with 25 oars on each side, while in the late Roman Republic, it was equipped with two banks of oars (a bireme), remaining faster, lighter, and more agile than triremes.
Maritime Southeast Asian warships:
Djong were sailing warships armed with up to a hundred cannons.
Kakap were small warships used in Maritime Southeast Asia.
Kelulus were Nusantaran warships used as troop transports and raiding vessels.
Lancaran were a type of galley warships armed with Cetbang cannons.
History and evolution of warships
First warships
The first evidence of ships being used for warfare comes from Ancient Egypt, specifically the northern Nile River most likely to defend against Mediterranean peoples. The galley warship most likely originated in Crete an idea which was soon copied and popularized by the Phoenicians. In the time of Mesopotamia, Ancient Persia, Phoenicia, Ancient Greece and the Ancient Rome, warships were always galleys (such as biremes, triremes and quinqueremes): long, narrow vessels powered by banks of oarsmen and designed to ram and sink enemy vessels, or to engage them bow-first and follow up with boarding parties. The development of catapults in the 4th century BC and the subsequent refinement of this technology enabled the first fleets of siege engine - equipped warships by the Hellenistic age. During late antiquity, ramming fell out of use and the galley tactics against other ships used during the Middle Ages until the late 16th century focused on boarding.
The Age of Sail
Naval artillery was redeveloped in the 14th century, but cannon did not become common at sea until the guns were capable of being reloaded quickly enough to be reused in the same battle. The size of a ship required to carry a large number of cannons made oar-based propulsion impossible, and warships came to rely primarily on sails. The sailing man-of-war emerged during the 16th century.
By the middle of the 17th century, warships were carrying increasing numbers of cannons on their broadsides and tactics evolved to bring each ship's firepower to bear in a line of battle. The man-of-war now evolved into the ship of the line. In the 18th century, the frigate and sloop-of-war – too small to stand in the line of battle – evolved to escort convoy trade, scout for enemy ships and blockade enemy coasts.
Steel, steam and shellfire
During the 19th century a revolution took place in the means of marine propulsion, naval armament and construction of warships. Marine steam engines were introduced, at first as an auxiliary force, in the second quarter of the 19th century. The Crimean War gave a great stimulus to the development of guns. The introduction of explosive shells soon led to the introduction of iron, and later steel, naval armour for the sides and decks of larger warships. The first ironclad warships, the French and British , made wooden vessels obsolete. Metal soon entirely replaced wood as the main material for warship construction.
From the 1850s, the sailing ships of the line were replaced by steam-powered battleships, while the sailing frigates were replaced by steam-powered cruisers.
The armament of warships also changed with the invention of the rotating barbettes and turrets, which allowed the guns to be aimed independently of the direction of the ship and allowed a smaller number of larger guns to be carried.
The final innovation during the 19th century was the development of the torpedo and development of the torpedo boat. Small, fast torpedo boats seemed to offer an alternative to building expensive fleets of battleships.
Pre-dreadnought era
Pre-dreadnought battleships were sea-going battleships built between the mid- to late- 1880s and 1905, before the launch of HMS Dreadnought in 1906. The pre-dreadnought ships replaced the ironclad battleships of the 1870s and 1880s. Built from steel, protected by case-hardened steel armour, and powered by coal-fired triple-expansion steam engines, pre-dreadnought battleships carried a main battery of very heavy guns in fully-enclosed rotating turrets supported by one or more secondary batteries of lighter weapons. The role of corvettes, sloops and frigates were taken by new types of ships like destroyers, protected cruisers and armoured cruisers.
Since 1906
The dreadnought era
Another revolution in capital warship design began shortly after the start of the 20th century, when Britain launched the Royal Navy's all-big-gun battleship in 1906. Powered by steam turbines, it was bigger, faster and more heavily gunned than any existing battleships, which it immediately rendered obsolete. It was rapidly followed by similar ships in other countries. The Royal Navy also developed the first battlecruisers. Mounting the same heavy guns as the dreadnoughts on an even larger hull, battlecruisers sacrificed armour protection for speed. Battlecruisers were faster and more powerful than all existing cruisers, but much more vulnerable to shellfire than contemporary battleships. The torpedo-boat destroyer was developed at the same time as the dreadnoughts. Bigger, faster and more heavily gunned than the torpedo boat, the destroyer evolved to protect the capital ships from the menace of the torpedo boat.
At this time, Britain also introduced the use of fuel oil to power steam warships, instead of coal. Oil produced twice as much power per unit weight as coal, and was much easier to handle. Tests were conducted by the Royal Navy in 1904 involving the torpedo-boat destroyer , the first warship powered solely by fuel oil. These proved its superiority, and all warships procured for the Royal Navy from 1912 were designed to burn fuel oil.
Obsolescence of battleships
During the lead-up to the Second World War, Germany and the United Kingdom once again emerged as the two dominant Atlantic sea powers. The German navy, under the Treaty of Versailles, was limited to only a few minor surface ships. But the clever use of deceptive terminology, such as Panzerschiffe deceived the British and French commands. They were surprised when ships such as , , and raided Allied supply lines. The greatest threat however, was the introduction of the Kriegsmarine's largest vessels, and . Bismarck was heavily damaged and sunk/scuttled after a series of sea battles in the north Atlantic in 1941, while Tirpitz was destroyed by the Royal Air Force in 1944. The British Royal Navy gained dominance of the European theatre by 1943.
The Second World War brought massive changes in the design and role of several types of warships. For the first time, the aircraft carrier became the clear choice to serve as the main capital ship within a naval task force. World War II was the only war in history in which battles occurred between groups of carriers. World War II saw the first use of radar in combat. It brought the first naval battle in which the ships of both sides never engaged in direct combat, instead sending aircraft to make the attacks, as in the Battle of the Coral Sea.
Cold War-era
Modern warships are generally divided into seven main categories, which are: aircraft carriers, cruisers, destroyers, frigates, corvettes, submarines, and amphibious warfare ships. Battleships comprise an eighth category, but are not in current service with any navy in the world. Only the deactivated American s still exist as potential combatants, and battleships in general are unlikely to re-emerge as a ship class without redefinition. The destroyer is generally regarded as the dominant surface-combat vessel of most modern blue-water navies. However, the once distinct roles and appearances of cruisers, destroyers, frigates, and corvettes have blurred. Most vessels have come to be armed with a mix of anti-surface, anti-submarine and anti-aircraft weapons. Class designations no longer reliably indicate a displacement hierarchy, and the size of all vessel types has grown beyond the definitions used earlier in the 20th century. Another key difference between older and modern vessels is that all modern warships are "soft", without the thick armor and bulging anti-torpedo protection of World War II and older designs.
Most navies also include many types of support and auxiliary vessels, such as minesweepers, patrol boats and offshore patrol vessels.
By 1982 the United Nations Convention on the Law of the Sea (UNCLOS) treaty negotiations had produced a legal definition of what was then generally accepted as a late-twentieth century warship. The UNCLOS definition was : "A warship means a ship belonging to the armed forces of a State bearing the external marks distinguishing such ships of its nationality, under the command of an officer duly commissioned by the government of the State and whose name appears in the appropriate service list or its equivalent, and manned by a crew which is under regular armed forces discipline."
Development of the submarine
The first practical submarines were developed in the late 19th century, but it was only after the development of the torpedo that submarines became truly dangerous (and hence useful). By the end of the First World War submarines had proved their potential. During the Second World War Nazi Germany's fleet of U-boats (submarines) almost starved Britain into submission and inflicted huge losses on US coastal shipping. The success of submarines led to the development of new anti-submarine convoy escorts during the First and Second World Wars, such as the destroyer escort. Confusingly, many of these new types adopted the names of the smaller warships from the age of sail, such as corvette, sloop and frigate.
Development of the aircraft carrier
A seaplane tender is a ship that supports the operation of seaplanes. Some of these vessels, known as seaplane carriers, could not only carry seaplanes but also provided all the facilities needed for their operation; these ships are regarded by some as the first aircraft carriers and appeared just before the First World War.
A major shift in naval warfare occurred with the introduction of the aircraft carrier. First at Taranto and then at Pearl Harbor, the aircraft carrier demonstrated its ability to strike decisively at enemy ships out of sight and range of surface vessels. By the end of the Second World War, the carrier had become the dominant warship.
Development of the amphibious assault ship
was a ship of the Imperial Japanese Army during World War II. She was the world's first purpose-built landing craft carrier ship, and a pioneer of modern-day amphibious assault ships. During some of her operations, she was known to have used at least four cover names, R1, GL, MT, and Ryujo Maru.
An amphibious warfare ship is an amphibious vehicle warship employed to land and support ground forces, such as marines, on enemy territory during an amphibious assault.
Specialized shipping can be divided into two types, most crudely described as ships and craft. In general, the ships carry the troops from the port of embarkation to the drop point for the assault and the craft carry the troops from the ship to the shore. Amphibious assaults taking place over short distances can also involve the shore-to-shore technique, where landing craft go directly from the port of embarkation to the assault point. Amphibious assault ships have a well deck with landing craft which can carry tanks and other armoured fighting vehicles and also have a deck like a helicopter carrier for helicopters and V/STOL aircraft.
| Technology | Naval warfare | null |
263964 | https://en.wikipedia.org/wiki/Osteoglossiformes | Osteoglossiformes | Osteoglossiformes (Greek: "bony tongues") is a relatively primitive order of ray-finned fish that contains two sub-orders, the Osteoglossoidei and the Notopteroidei. All of at least 245 living species inhabit freshwater. They are found in South America, Africa, Australia and southern Asia, having first evolved in Gondwana before that continent broke up. In 2008, several new species of marine osteoglossiforms were described from the Danish Eocene Fur Formation, dramatically increasing the diversity of this group. This implies that the Osteoglossomorpha is not a primary freshwater fish group with the osteoglossiforms having a typical Gondwana distribution.
The Gymnarchidae (the only species being Gymnarchus niloticus, the African knifefish) and the Mormyridae are weakly electric fish able to sense their prey using electric fields.
The mooneyes (Hiodontidae) are often classified here, but may also be placed in a separate order, Hiodontiformes.
Members of the order are notable for having toothed or bony tongues, and for having the forward part of the gastrointestinal tract pass to the left of the oesophagus and stomach (for all other fish it passes to the right). In other respects, osteoglossiform fishes vary considerably in size and form; the smallest is Pollimyrus castelnaui, at just long, while the largest, the arapaima (Arapaima gigas), reaches as much as .
Phylogeny
Phylogeny based on the following works:
Families
Eschmeyer's Catalog of Fishes does not recognise suborders within the order Osetglossiformes and recognises the following families:
Family Pantodontidae Peters, 1876 (freshwater butterfly fishes)
Family Osteoglossidae Bonaparte, 1845 (bonytongues)
Family Arapaimidae Bonaparte, 1846 (arapaimas)
Family Notopteridae Bleeker, 1851 (featherfin knifefishes or featherbacks)
Subfamily Notopterinae Bleeker, 1851 (Asiatic featherbacks)
Subfamily Xenomystinae Greenwood, 1963 (African knifefishes)
Family Mormyridae Bonaparte, 1831 (elephantfishes)
Family Gymnarchidae Bleeker, 1859 (abas)
| Biology and health sciences | Osteoglossiformes | Animals |
263970 | https://en.wikipedia.org/wiki/Volcanic%20cone | Volcanic cone | Volcanic cones are among the simplest volcanic landforms. They are built by ejecta from a volcanic vent, piling up around the vent in the shape of a cone with a central crater. Volcanic cones are of different types, depending upon the nature and size of the fragments ejected during the eruption. Types of volcanic cones include stratocones, spatter cones, tuff cones, and cinder cones.
Stratocone
Stratocones are large cone-shaped volcanoes made up of lava flows, explosively erupted pyroclastic rocks, and igneous intrusives that are typically centered around a cylindrical vent. Unlike shield volcanoes, they are characterized by a steep profile and periodic, often alternating, explosive eruptions and effusive eruptions. Some have collapsed craters called calderas. The central core of a stratocone is commonly dominated by a central core of intrusive rocks that range from around to over several kilometers in diameter. This central core is surrounded by multiple generations of lava flows, many of which are brecciated, and a wide range of pyroclastic rocks and reworked volcanic debris. The typical stratocone is an andesitic to dacitic volcano that is associated with subduction zones. They are also known as either stratified volcano, composite cone, bedded volcano, cone of mixed type or Vesuvian-type volcano.
Spatter cone
A spatter cone is a low, steep-sided hill or mound that consists of welded lava fragments, called spatter, which has formed around a lava fountain issuing from a central vent. Typically, spatter cones are about high. In case of a linear fissure, lava fountaining will create broad embankments of spatter, called spatter ramparts, along both sides of the fissure. Spatter cones are more circular and cone shaped, while spatter ramparts are linear wall-like features.
Spatter cones and spatter ramparts are typically formed by lava fountaining associated with mafic, highly fluid lavas, such as those erupted in the Hawaiian Islands. As blobs of molten lava, spatter, are erupted into the air by a lava fountain, they can lack the time needed to cool completely before hitting the ground. Consequently, the spatter are not fully solid, like taffy, as they land and they bind to the underlying spatter as both often slowly ooze down the side of the cone. As a result, the spatter builds up a cone that is composed of spatter either agglutinated or welded to each other.
Tuff cones
A tuff cone, sometimes called an ash cone, is a small monogenetic volcanic cone produced by phreatic (hydrovolcanic) explosions directly associated with magma brought to the surface through a conduit from a deep-seated magma reservoir. They are characterized by high rims that have a maximum relief of above the crater floor and steep slopes that are greater than 25 degrees. They typically have a rim to rim diameter of . A tuff cone consists typically of thick-bedded pyroclastic flow and surge deposits created by eruption-fed density currents and bomb-scoria beds derived from fallout from its eruption column. The tuffs composing a tuff cone have commonly been altered, palagonitized, by either its interaction with groundwater or when it was deposited warm and wet. The pyroclastic deposits of tuff cones differ from the pyroclastic deposits of spatter cones by their lack or paucity of lava spatter, smaller grain-size, and excellent bedding. Typically, but not always, tuff cones lack associated lava flows.
A tuff ring is a related type of small monogenetic volcano that is also produced by phreatic (hydrovolcanic) explosions directly associated with magma brought to the surface through a conduit from a deep-seated magma reservoir. They are characterized by rims that have a low, broad topographic profiles and gentle topographic slopes that are 25 degrees or less. The maximum thickness of the pyroclastic debris comprising the rim of a typical tuff ring is generally thin, less than to thick. The pyroclastic materials that comprise their rim consist primarily of relatively fresh and unaltered, distinctly and thin-bedded volcanic surge and air fall deposits. Their rims also can contain variable amounts of local country rock (bedrock) blasted out of their crater. In contrast to tuff cones, the crater of a tuff ring generally has been excavated below the existing ground surface. As a result, water commonly fills a tuff ring's crater to form a lake once eruptions cease.
Both tuff cones and their associated tuff rings were created by explosive eruptions from a vent where the magma is interacting with either groundwater or a shallow body of water as found within a lake or sea. The interaction between the magma, expanding steam, and volcanic gases resulted in the production and ejection of fine-grained pyroclastic debris called ash with the consistency of flour. The volcanic ash comprising a tuff cone accumulated either as fallout from eruption columns, from low-density volcanic surges and pyroclastic flows, or combination of these. Tuff cones are typically associated with volcanic eruptions within shallow bodies of water and tuff rings are associated with eruptions within either water saturated sediments and bedrock or permafrost.
Next to spatter (scoria) cones, tuff cones and their associated tuff rings are among the most common types of volcanoes on Earth. An example of a tuff cone is Diamond Head at Waikīkī in Hawaii. Clusters of pitted cones observed in the Nephentes/Amenthes region of Mars at the southern margin of the ancient Utopia impact basin are currently interpreted as being tuff cones and rings.
Cinder cone
Cinder cones, also known as scoria cones and less commonly scoria mounds, are small, steep-sided volcanic cones built of loose pyroclastic fragments, such as either volcanic clinkers, cinders, volcanic ash, or scoria. They consist of loose pyroclastic debris formed by explosive eruptions or lava fountains from a single, typically cylindrical, vent. As the gas-charged lava is blown violently into the air, it breaks into small fragments that solidify and fall as either cinders, clinkers, or scoria around the vent to form a cone that often is noticeably symmetrical; with slopes between 30 and 40°; and a nearly circular ground plan. Most cinder cones have a bowl-shaped crater at the summit. The basal diameters of cinder cones average about and range from . The diameter of their craters ranges between . Cinder cones rarely rise more than or so above their surroundings.
Cinder cones most commonly occur as isolated cones in large basaltic volcanic fields. They also occur in nested clusters in association with complex tuff ring and maar complexes. Finally, they are also common as parasitic and monogenetic cones on complex shield and stratovolcanoes. Globally, cinder cones are the most typical volcanic landform found within continental intraplate volcanic fields and also occur in some subduction zone settings as well. Parícutin, the Mexican cinder cone which was born in a cornfield on February 20, 1943, and Sunset Crater in Northern Arizona in the US Southwest are classic examples of cinder cones, as are ancient volcanic cones found in New Mexico's Petroglyph National Monument. Cone-shaped hills observed in satellite imagery of the calderas and volcanic cones of Ulysses Patera, Ulysses Colles and Hydraotes Chaos are argued to be cinder cones.
Cinder cones typically only erupt once like Parícutin. As a result, they are considered to be monogenetic volcanoes and most of them form monogenetic volcanic fields. Cinder cones are typically active for very brief periods of time before becoming inactive. Their eruptions range in duration from a few days to a few years. Of observed cinder cone eruptions, 50% have lasted for less than 30 days, and 95% stopped within one year. In case of Parícutin, its eruption lasted for nine years from 1943 to 1952. Rarely do they erupt either two, three, or more times. Later eruptions typically produce new cones within a volcanic field at separation distances of a few kilometers and separate by periods of 100 to 1,000 years. Within a volcanic field, eruptions can occur over a period of a million years. Once eruptions cease, being unconsolidated, cinder cones tend to erode rapidly unless further eruptions occur.
Rootless cones
Rootless cones, also called pseudocraters, are volcanic cones that are not directly associated with a conduit that brought magma to the surface from a deep-seated magma reservoir. Generally, three types of rootless cones, littoral cones, explosion craters, and hornitos are recognized. Littoral cones and explosion craters are the result of mild explosions that were generated locally by the interaction of either hot lava or pyroclastic flows with water. Littoral cones typically form on the surface of a basaltic lava flow where it has entered into a body of water, usually a sea or ocean. Explosion craters form where either hot lava or pyroclastic flows have covered either marshy ground or water-saturated ground of some sort. Hornitos are rootless cones that are composed of welded lava fragments and were formed on the surface of basaltic lava flows by the escape of gas and clots of molten lava through cracks or other openings in the crust of a lava flow.
| Physical sciences | Volcanic landforms | Earth science |
264025 | https://en.wikipedia.org/wiki/Curassow | Curassow | Curassows are one of the three major groups of cracid birds. They comprise the largest-bodied species of the cracid family. Three of the four genera are restricted to tropical South America; a single species of Crax ranges north to Mexico. They form a distinct clade which is usually classified as the subfamily Cracinae.
Evolution
In line with the other 3 main lineages of cracids (chachalacas, true guans, and the horned guan), mt and nDNA sequence data indicates that the curassows diverged from their closest living relatives (probably the guans) at some time during the Oligocene, or c.35–20 mya (Pereira et al. 2002). This data must be considered preliminary until corroborated by material (e.g. fossil) evidence however.
What appears certain from analysis of the molecular data, calibrated against geological events that would have induced speciation, is that there are 2 major lineages of curassows: one containing only Crax, and another made up of Mitu and Pauxi. The position of the peculiar nocturnal curassow is not well resolved; it might be closer to the latter, but in any case, it diverged around the same time as the split between the two major lineages. All curassow genera appear to have diverged, in fact, during the Tortonian (early Late Miocene): the initial split took place some 10–9 mya, and Pauxi diverged from Mitu some 8–7.4 mya (but see genus article).(Pereira & Baker 2004)
Unlike the other cracids, biogeography and phylogeny indicate that the extant lineages of curassows probably originated in the lowlands of the western/northwestern Amazonas basin, most likely in the general area where the borders of Brazil, Peru and Colombia meet. In the two larger genera, vicariant speciation seems to have played a major role.(Pereira et al. 2002, Pereira & Baker 2004)
| Biology and health sciences | Galliformes | Animals |
264210 | https://en.wikipedia.org/wiki/Zero%20of%20a%20function | Zero of a function | In mathematics, a zero (also sometimes called a root) of a real-, complex-, or generally vector-valued function , is a member of the domain of such that vanishes at ; that is, the function attains the value of 0 at , or equivalently, is a solution to the equation . A "zero" of a function is thus an input value that produces an output of 0.
A root of a polynomial is a zero of the corresponding polynomial function. The fundamental theorem of algebra shows that any non-zero polynomial has a number of roots at most equal to its degree, and that the number of roots and the degree are equal when one considers the complex roots (or more generally, the roots in an algebraically closed extension) counted with their multiplicities. For example, the polynomial of degree two, defined by has the two roots (or zeros) that are 2 and 3.
If the function maps real numbers to real numbers, then its zeros are the -coordinates of the points where its graph meets the x-axis. An alternative name for such a point in this context is an -intercept.
Solution of an equation
Every equation in the unknown may be rewritten as
by regrouping all the terms in the left-hand side. It follows that the solutions of such an equation are exactly the zeros of the function . In other words, a "zero of a function" is precisely a "solution of the equation obtained by equating the function to 0", and the study of zeros of functions is exactly the same as the study of solutions of equations.
Polynomial roots
Every real polynomial of odd degree has an odd number of real roots (counting multiplicities); likewise, a real polynomial of even degree must have an even number of real roots. Consequently, real odd polynomials must have at least one real root (because the smallest odd whole number is 1), whereas even polynomials may have none. This principle can be proven by reference to the intermediate value theorem: since polynomial functions are continuous, the function value must cross zero, in the process of changing from negative to positive or vice versa (which always happens for odd functions).
Fundamental theorem of algebra
The fundamental theorem of algebra states that every polynomial of degree has complex roots, counted with their multiplicities. The non-real roots of polynomials with real coefficients come in conjugate pairs. Vieta's formulas relate the coefficients of a polynomial to sums and products of its roots.
Computing roots
There are many methods for computing accurate approximations of roots of functions, the best being Newton's method, see Root-finding algorithm.
For polynomials, there are specialized algorithms that are more efficient and may provide all roots or all real roots; see Polynomial root-finding and Real-root isolation.
Some polynomial, including all those of degree no greater than 4, can have all their roots expressed algebraically in terms of their coefficients; see Solution in radicals.
Zero set
In various areas of mathematics, the zero set of a function is the set of all its zeros. More precisely, if is a real-valued function (or, more generally, a function taking values in some additive group), its zero set is , the inverse image of in .
Under the same hypothesis on the codomain of the function, a level set of a function is the zero set of the function for some in the codomain of
The zero set of a linear map is also known as its kernel.
The cozero set of the function is the complement of the zero set of (i.e., the subset of on which is nonzero).
Applications
In algebraic geometry, the first definition of an algebraic variety is through zero sets. Specifically, an affine algebraic set is the intersection of the zero sets of several polynomials, in a polynomial ring over a field. In this context, a zero set is sometimes called a zero locus.
In analysis and geometry, any closed subset of is the zero set of a smooth function defined on all of . This extends to any smooth manifold as a corollary of paracompactness.
In differential geometry, zero sets are frequently used to define manifolds. An important special case is the case that is a smooth function from to . If zero is a regular value of , then the zero set of is a smooth manifold of dimension by the regular value theorem.
For example, the unit -sphere in is the zero set of the real-valued function .
| Mathematics | Functions: General | null |
264213 | https://en.wikipedia.org/wiki/Sheremetyevo%20International%20Airport | Sheremetyevo International Airport | Sheremetyevo Alexander S. Pushkin International Airport () is one of four international airports that serve the city of Moscow. It is the busiest airport in Russia, as well as the 11th-busiest airport in Europe. Originally built as a military airbase, Sheremetyevo was converted into a civilian airport in 1959. The airport was originally named after a nearby village, and a 2019 contest extended the name to include the name of the Russian poet Alexander Pushkin.
The airport comprises six terminals: four international terminals (one under construction), one domestic terminal, and one private aviation terminal. It is located northwest of central Moscow, between the towns of Lobnya and Khimki in Moscow Oblast.
In 2019, the airport handled about 49.9 million passengers. Sheremetyevo serves as the main hub for Russian flag carrier Aeroflot as well as its subsidiaries Rossiya Airlines and Pobeda, for Nordwind Airlines and its subsidiary Ikar, and for Smartavia.
History
Soviet era
The airport was initially built as a military airfield called Sheremetyevsky (), named after a village of the same name, as well as the nearby railway station of the same name. The decree for the construction of the Central Airdrome of the Air Force near the settlement of Chashnikovo on the outskirts of Moscow was issued on 1 September 1953 by the Council of Ministers of the Soviet Union. The airport became operational on 7 November 1957 to celebrate the 40th anniversary of the October Revolution.
In August 1959, the Council of Ministers made a decree to terminate the airbase's use for military purposes, where it would be handed over to the Principal Directorate of the Civil Air Fleet to be converted into a civilian airport. Sheremetyevo's civilian purposes started on 11 August 1959 when a Tupolev Tu-104B landed at the airport from Leningrad.
The first international flight took place on 1 June 1960 to Berlin Schönefeld Airport using an Ilyushin Il-18. Sheremetyevo was officially opened on the day after, where a two-story terminal occupying was commissioned. On 3 September 1964, the Sheremetyevo-1 terminal was opened. Of that year, 18 foreign airlines had regular flights to Sheremetyevo, with up to 10 different types of aircraft involved. By the end of 1964, Sheremetyevo handled 822,000 passengers and 23,000 tons of mail and cargo, including 245,000 passengers and 12,000 tons of cargo that were transported internationally. Soon, by the end of 1965, a majority of international flights to the USSR was achieved through Sheremetyevo thanks to Aeroflot's air traffic agreements with 47 countries.
In the early 1970s, a second runway was constructed at Sheremetyevo, with the first airliner to land being an Ilyushin Il-62. In preparation for the 1980 Summer Olympics, construction of a second terminal for Sheremetyevo, Sheremetyevo-2, was approved by the Ministry of Civil Aviation in early 1976. Construction of Sheremetyevo-2 started on 17 November 1977.
On 1 January 1980, Sheremetyevo-2 was put into operation, with a capacity to serve an annual 6 million passengers, or 2,100 passengers per hour. Despite this, its official opening ceremony was held much later, on 6 May 1980. During the Olympics, Sheremetyevo served more than 460,000 international passengers.
Contemporary era
On 11 November 1991, Sheremetyevo International Airport received its legal status as a state-owned enterprise, amidst the dissolution of the Soviet Union. On 9 July 1996, Sheremetyevo became an open joint-stock company. In 1997, the airport renovated one of its runways with a thick concrete surface.
In the early 2000s, Sheremetyevo saw growing competition from the rapidly expanding Moscow Domodedovo Airport, which was more modern and convenient to access, and the neighbouring Vnukovo Airport. Sheremetyevo saw 24 of its airlines, notably domestic airlines such as Sibir, KrasAir, Transaero, Pulkovo Airlines, and UTAir, as well as international airlines Air Malta, Adria Airlines, Swiss, British Airways, and Emirates, move their services to Domodedovo. As a result, Aeroflot pushed for a third terminal for the airport, Sheremetyevo-3, to increase the airport's passenger capacity as well as be able to fulfill its requirements to join Skyteam.
In the late 2000s, Sheremetyevo oversaw rapid planning and expansion of the airport. On 12 March 2007, the airport opened Terminal C to maximise the airport's international passenger capacity. On 5 March 2008, the airport renovated its second runway to receive all types of aircraft, including the Airbus A380 and the Boeing 787 Dreamliner. An Aeroexpress line was constructed between Sheremetyevo and Savyolovsky Railway Station on 10 June 2008, quickening traveling time from the airport to central Moscow in 30 minutes. In January 2009, Sheremetyevo finalised a master plan where it would increase passenger capacity to an annual 64 million per year and build a second airfield with a third runway. On 15 November 2009, construction of Terminal D was completed, with a total surface area of , an annual capacity of 12 million passengers, and operation being putting forth in the beginning of next year. Sheremetyevo-2 was renamed Terminal F on 25 December 2009 with terminal identification using international (Latin) lettering.
Expansion of Sheremetyevo continued into 2010. Sheremetyevo-1 was renamed Terminal B on 28 March. Terminal E was opened on 30 April, connecting Terminal D and Terminal F and increasing the airport's capacity to 35 million passengers per year. In June, construction started for Terminal A, a private aviation terminal. In July, a walkway opened between Terminals D, E, F, and the Aeroexpress railway terminal on the public access side. In November, a walkway opened between Terminals D, E, and F on the security side. Both of have simplified transfer between transit flights. Ultimately, after the northern the recent construction work, the airport now has the capacity to receive more than 40 million passengers annually.
On 28 March 2011, a separate airfield that would serve as Sheremetyevo's third runway was approved. On 13 December 2011, the Federal Agency for State Property Management approved an agreement that merged the airport operators OAO Terminal (operator of Terminal D) and OJSC Sheremetyevo, consolidating control of the airport under one entity. On 26 December 2011, a new area control centre (ACC) was opened for Sheremetyevo, consolidating operations of the airport's different control centres to increase efficiency. The situational centre was also created as part of the ACC for joint work of top-managers, heads of state bodies, and partners of Sheremetyevo to resolve emergencies.
Continued expansion
On 30 December 2013, TPS Avia successfully won a competitive tender to develop Sheremetyevo International Airport's northern area, including a new passenger terminal, a new freight terminal, a refuelling area and a tunnel linking the passenger terminal to three other terminals.
Terminal B, previously Sheremetyevo-1, was demolished in August 2015 to be reconstructed as a newer and more modern terminal, which began in October 2015. By the end of 2015, Sheremetyevo surpassed its competitor Domodedovo as Russia's busiest airport, serving 31.28 million passengers, compared to Domodedovo's 30.05 million. This trend continued in 2016, where Sheremetyevo saw growth while Vnukovo and Domodedovo showed losses in passengers. A growing number of airlines launched new operations to Sheremetyevo, such as Tianjin Airlines, Tunisair, Nouvelair, and Air Malta, which back in the 2000s moved its operation to Domodedovo.
In February 2016, TPS Avia combined its assets with Sheremetyevo Airport and committed to invest US$840 million to upgrade and expand the airport's infrastructure – as a result TPS Avia secured a 68% stake in Sheremetyevo Airport. Part of the plan includes demolishing Terminal C for a newer reconstruction of the terminal, which came to effect on 1 April 2017. Sheremetyevo International Airport was the official airport of the 2018 FIFA World Cup. Sheremetyevo completed re-construction of its first northern terminal, Terminal B, in May 2018, to handle more passengers for the tournament. In 2018, the Airport reported revenues of €194.9 million, a 6% increase year over year. Profit increased 7.4% year over year. These increases are attributed in part to increased air traffic due to the 2018 FIFA World Cup.
In late 2018, SVO enacted a series of changes to its flight traffic. Aeroflot subsidiary Rossiya Airlines announced the transfer of its flights from Vnukovo to Sheremetyevo starting 28 October 2018. British Airways also launched direct flights from London Heathrow to Sheremetyevo on the same day. Syria-based Cham Wings Airlines began direct flights from Damascus to SVO in November 2018 as well. In December 2018, following the results of the Great Names of Russia contest, Sheremetyevo was named after the great Russian poet Alexander Pushkin. The ceremony took place on 5 June 2019, which was the 220th anniversary of Pushkin's birth year. The airport is now officially named Sheremetyevo Alexander S. Pushkin International Airport.
In 2019, the Russian Federal Security Service (FSB) began testing an automated passport control system at SVO. This system relies on biometric data and foreign passport recognition to allow Russian passengers to move through border control with fewer movement restrictions. If successful, the FSB may implement this system in other Russian airports.
Fraud
OCCRP reports that the airport has been used for laundering money. It purchased fuel from a broad network of middlemen between 2003 and 2008, which greatly increased the price. Court records show that just in 2006 and 2007, phantom corporations made at least $200 million in pointless markups. The scam cost the Russian government approximately 1 billion rubles ($40 million) in missing tax income. The cost of fuel increased, which also increased the cost of airline tickets for the general people.
Terminals
Sheremetyevo International Airport has four operating passenger terminals and one special terminal reserved for the use of private and business aviation. The airport's four passenger terminals are divided into two groups based on geographical location: the Northern Terminal Complex and the Southern Terminal Complex. The current terminal naming system was introduced in December 2009; previously, the terminals were numbered: Sheremetyevo-1 (now Terminal B), Sheremetyevo-2 (now Terminal F), and Sheremetevo-3 (now Terminal D).
Terminal A
Opened on 16 January 2012, Terminal A handles servicing of business and private aviation out of Sheremetyevo. The terminal occupies an area of and can carry an annual capacity of 75,000 passengers.
Northern terminals
Terminal B
Terminal B – originally named Sheremetyevo-1 – has two iterations.
The first iteration was constructed and opened on 3 September 1964. The terminal, as Sheremetyevo-1, was known for its "flying-saucer"-like design, and was nicknamed "shot glass" by locals. Being long and wide, as well as having a volume exceeding , the terminal can hold up to 800 people per hour. Formerly serving international flights, Sheremetyevo-1 would transition to serving domestic flights. Along with other Sheremetyevo terminals that underwent Latin lettering conventions, Sheremetyevo-1 was renamed Terminal B on 28 March 2010. Terminal B was then demolished in August 2015 to be reconstructed as a larger and more modern terminal which began in October 2015.
The new terminal B commenced its operations on 3 May 2018, with the Aeroflot's flight to Saratov. All airlines that have domestic flights from Sheremetyevo and some flights of Aeroflot began shifting to Terminal B from Terminal D. Compared to the previous terminal B, that was demolished, new terminal will have an increased passenger capacity of 20 million passengers and will serve domestic flights only. As of November 2018, Aeroflot has consolidated all of its domestic services at Terminal B, with the exception of flights to far eastern destinations in Vladivostok, Khabarovsk and Petropavlovsk-Kamchatsky. Flights to the eastern Russian shore and some short-haul (including all domestic flights served by widebodies) continue out of SVO's Terminal D.
The terminal is connected by an interterminal underground passage with Sheremetyevo's southern terminals and the Aeroexpress railway station.
Terminal C
On 12 March 2007, Sheremetyevo opened the former Terminal C for the servicing of international charter flights to maximize location convenience for all areas in the airport. Located adjacent to the former Terminal B, Terminal C served from 5 to 6 million passengers. The role of Terminal C diminished as passengers for international flights for the airport were distributed among Terminal D and Terminal E. As part of Sheremetyevo's long-term redevelopment plan, Terminal C was closed on 1 April 2017 to be demolished for construction of a newer terminal.
Integrated with the now-reconstructed domestic Terminal B, the new Terminal C was designed to serve up to 20 million passengers.
The first section of the new Terminal C opened on 17 January 2020, with a planned capacity of 20 million passengers. It is called Terminal C1, and some international flights were transferred to that new terminal. Another part called Terminal C2 is scheduled to be opened in 2026, and will add another 10 million passengers capacity.
Southern terminals
Terminal D
Terminal D, opened in November 2009, is adjacent to Terminal F. The building is a hub for Aeroflot and its SkyTeam partners, with capacity for 12 million passengers per year. Aeroflot had been trying to implement the project of a new terminal (Sheremetyevo-3) since January 2001. However, construction only began in 2005, with commissioning of the complex finally taking place on 15 November 2009. The acquisition of its own terminal was a condition of Aeroflot's entry into the SkyTeam airline alliance, thus necessitating the construction. The main contractor for the build was a Turkish company Enka. Terminal D has 22 jetways and 11 remote stands. On 15 November 2009 at 9:15 a.m., the first flight from Terminal D (the new official name of Sheremetyevo-3) departed for the southern resort city of Sochi. Despite this, Aeroflot took a number of months (due to unexpected administrative delays) to transfer all of its international flights from Terminal F to D (a full transfer was originally planned for February 2010). Whilst previously Terminal D had remained a separate legal entity from the rest of Sheremetyevo Airport, in spring 2012, it became an integrated unit of "Sheremetyevo International Airport" JSC. As part of the deal, Aeroflot, VEB Bank, and VTB Bank, all of which had invested in the construction of Terminal D, became part shareholders in the airport as a whole. The basis for the architectural and artistic image of Terminal D is that of a giant swan with outstretched wings.
There is an official multi-story parking at Terminal D connected with the main building by means of a pedestrian bridge. The parking size is about 4100 lots, however it has a relatively dense layout.
Between August 2015 and May 2018, Terminal D used to be the only terminal at Sheremetyevo that was able to serve domestic flights. Even since new Terminal B was opened and commenced its services, Terminal D continues to operate non-Aeroflot domestic flights.
On 28 October 2018, Terminal D started handling all of Rossiya Airlines' Moscow-originating domestic flights and its international service to Indonesia.
On 15 March 2022, the Terminal D was closed caused by dramatic passenger traffic decrease. On terminal was reopened. As of October 2024 it serve flights operated by Pobeda and Smartavia.
Terminal E
Terminal E opened in 2010 as a capacity expansion project, connecting terminals D and F. The terminal's construction has allowed for the development of terminals D and F, as well as the railway station, into a single south terminal complex. The terminals of this complex are connected by a number of pedestrian walkways with travelators, thus allowing for passengers to move freely between its constituent facilities. In December 2010, a new chapel dedicated to St. Nicholas opened on the second floor of Terminal E. The terminal was used for international flights, primarily by Aeroflot and its SkyTeam partners. Terminal E has 8 jetway equipped gates. The V-Express Transit Hotel between security/passport check-ins provided short-term accommodations for passengers changing planes without having to present a visa for entering Russia. The hotel drew international attention in June 2013 when Edward Snowden checked into the hotel while seeking asylum.
In March 2020, Terminal E was closed due to a decrease in passenger flow and due to COVID-19 in Russia.
Terminal F
Opened on 6 May 1980 for Moscow's Summer Olympics, Terminal F, previously Sheremetyevo-2, has 15 jetways and 21 remote aircraft stands. The terminal was designed to service 6 million passengers per year. Until the completion of the original Terminal C, it was the only terminal that serviced international flights. The design is a larger version of the one of Hannover–Langenhagen Airport by the same architects and constructed by Rüterbau, a company located in Hanover. All materials, except the bricks which came from Poland, and every piece of equipment, was transported from Germany to Moscow by lorry. A major reconstruction of the terminal and its interior space was completed by late 2009. For the convenience of passengers, the departures lounge and duty free zone were thoroughly modernised, whilst a number of partition walls were removed to create extra retail and lounge space.
It was announced that terminal F will be re-constructed after the construction of terminal C is completed.
On 30 December 2021, at 0:00 by Moscow Time, the terminal F was closed for reconstruction.
Terminal G
In November 2019, it was announced that a new Terminal G will also be built. Construction is planned to begin in 2024-2025.
Airlines and destinations
Passenger
The following airlines serve regular scheduled and charter destinations at Sheremetyevo International Airport.
Cargo
Statistics
Ground transport
Rail
Aeroexpress, a subsidiary of Russian Railways operates a nonstop line, connecting the airport to Belorussky station in downtown Moscow. A one-way journey takes 35 minutes. The trains offer adjustable seats, luggage compartments, restrooms, electric outlets. Business-class coaches available.
The service started in November 2004, when express train connection was established from Savyolovsky station to Lobnya station, which is from the airport, with the remainder of the journey served by bus or taxi. On 10 June 2008, a rail terminal opened in front of Terminal F, with direct service from Savyolovsky station. A shuttle bus service ferried passengers to terminals B and C. From 28 August 2009, the line was extended to Belorussky station with plans to serve all three of Moscow's main airports from a single point of boarding, and service to Savyolovsky station terminated.
Interterminal underground
The airport's Automated Passenger Transportation System (APTS) connects the Terminal B and C with the Terminals D, E, F and the Aeroexpress railway station.
At the 1st floor of the Terminal B there is an entrance to Sheremetyevo 1 — the northern station. The entrance to Sheremetyevo 2 — the southern station — is at the passage between the terminals D and E.
The APTS is a part of the — a dual tunnel transportation system in the airport. One of the tunnels is dedicated to the transportation of people and featuring an automated people mover (APM). The other tunnel is used for automated baggage transportation.
Bus
Moscow can be reached by the municipal Mosgortrans bus lines: 817 to station Planernaya of Moscow Metro Tagansko-Krasnopresnenskaya Line (#7), 851 to station Rechnoy Vokzal of Zamoskvoretskaya Line (#2), departures every 10 minutes, travel time 33–55 minutes by schedule depending on the terminal served. At night time bus N1 () (departures every 30 minutes between 3am and 5:40am) connects the airport to Moscow's Leningradsky Avenue, downtown area and Leninsky Avenue. Travel time 30–90 minutes, fare is 57 rubles (as of February 2021).
Other buses serve the connections to the nearby cities: Lobnya (route 21), Zelenograd, Khimki (routes 43,62), Dolgoprudny.
Road
The main road leading to the airport—Leningradskoye Highway—has experienced large traffic jams. Since 23 December 2014, a toll road to the airport has been opened. It connects with MKAD near Dmitrovskoe Highway. Now it is possible to reach the airport in ten minutes, avoiding traffic jams.
Official airport taxis are available from taxi counters in arrivals. Prices to the city are fixed based on zones.
Accidents and incidents
On 26 September 1960, Austrian Airlines Flight 901 crashed short of the runway at Sheremetyevo Airport. Of the 37 people on board, 31 died.
On 28 November 1972, Japan Airlines Flight 446, a DC-8-62, crashed while in an initial climb on a route from Sheremetyevo International Airport to Haneda Airport. There were 14 crew members and 62 board the aircraft. A total of 9 crew and 52 passengers died, with a total of 61 of 76 occupants dead.
On 28 November 1976, Aeroflot Flight 2415, a Tupolev Tu-104 crashed shortly after takeoff as result of artificial horizon failure. All 67 passengers and six crew members died in the crash.
On 6 July 1982, Aeroflot Flight 411, an Ilyushin Il-62, crashed on takeoff; all 90 on board died.
On 22 July 2002, Pulkovo Aviation Enterprise Flight 9560, an Ilyushin Il-86, crashed on takeoff; 14 of the 16 occupants on board died.
On 3 June 2014, Ilyushin Il-96 RA-96010 of Aeroflot was damaged beyond economical repair in a fire whilst parked.
On 5 May 2019, Aeroflot Flight 1492, a Sukhoi Superjet 100, crash-landed and caught fire after returning to the airport due to an on-board malfunction shortly after takeoff, killing 41 of the 78 passengers and crew on board and injuring 11 others.
Awards and accolades
In 2018, Sheremetyevo International Airport was recognized for the best customer service in the busiest airports in Europe category by ACI's global Airport Service Quality (ASQ) program. In 2018, Sheremetyevo entered the list of the world's best airports – ACI Director General's Roll of Excellence. The Official Aviation Guide (OAG) ranked Sheremetyevo International Airport as the most punctual major airport (20 – 30 million departing seats) in the world for 2018, with an on-time performance of 87%.
In February 2019, SVO won an award for strengthening Russia's national security with its perimeter protection system. In February 2019, Sheremetyevo on top in on-time departure performance in the Major Airports category for February 2019, with 93.65% flights departed on time. In March 2019, Sheremetyevo International Airport was officially awarded a 5-star terminal rating from Skytrax, with Terminal B receiving the 5-star rating after a comprehensive audit.
In January 2020, Sheremetyevo International Airport has been named by the travel data and analytics expert Cirium as the world's most punctual airport in the annual On-Time Performance (OTP) review, with 95% of its flights departing on-time.
Sheremetyevo International Airport was recognized as the best airport for service quality in 2020 among airports with 2019 passenger traffic of more than 40 million by the Airports Council International's (ACI) global program for researching the level of service at airports Airport Service Quality (ASQ). At the end of 2020, Sheremetyevo topped the rating in the category of the largest airports in Europe for the third time. At the same time, this year Sheremetyevo was included in the list of the Voice of the Customer of the Airports Council International – the 140 most active airports in the implementation of the ASQ ACI program during the COVID-19 pandemic.
| Technology | Europe | null |
264501 | https://en.wikipedia.org/wiki/Vicia%20faba | Vicia faba | Vicia faba, commonly known as the broad bean, fava bean, or faba bean, is a species of vetch, a flowering plant in the pea and bean family Fabaceae. It is widely cultivated as a crop for human consumption, and also as a cover crop. Varieties with smaller, harder seeds that are fed to horses or other animals are called field bean, tic bean or tick bean. Horse bean, Vicia faba var. equina Pers., is a variety accepted as distinct. This legume is very common in Southern European, Northern European, East Asian, Latin American and North African cuisines.
Some people suffer from favism, a hemolytic response to the consumption of broad beans, a condition linked to a metabolism disorder known as G6PDD. Otherwise the beans, with the outer seed coat removed, can be eaten raw or cooked. In young plants, the outer seed coat can be eaten, and in very young plants, the seed pod can be eaten.
Description
Vicia faba is a stiffly erect, annual plant tall, with two to four stems that are square in cross-section. The leaves are long, pinnate with 2–7 leaflets, and glaucous (grey-green). Unlike most other vetches, the leaves do not have tendrils for climbing over other vegetation.
The flowers are long with five petals; the standard petals are white, the wing petals are white with a black spot (true black, not deep purple or blue as is the case in many "black" markings) and the keel petals are white. Crimson-flowered broad beans also exist, which were recently saved from extinction. The flowers have a strong sweet scent which is attractive to bees and other pollinators.
The fruit is a broad, leathery pod that is green, but matures to a dark blackish-brown, with a densely downy surface; the wild species has pods that are long and 1 cm diameter, but many modern cultivars developed for food use have pods long and 2–3 cm thick. Each bean pod contains 3–8 seeds. They are round to oval and have a 5–10 mm diameter in the wild plant, but are usually flattened and up to 20–25 mm long, 15 mm broad and 5–10 mm thick in food cultivars. V. faba has a diploid (2n) chromosome number of 12 (six homologous pairs). Five pairs are acrocentric chromosomes and one pair is metacentric.
Genome
The diploid genome of Vicia faba contains 13 GB of DNA, mostly obtained through amplification of retrotransposons and satellite repeats. The genome is one of the largest diploid field crops and contains a predicted 34,221 protein-coding genes.
Cultivation
Broad beans have a long tradition of cultivation in Old World agriculture, being among the most ancient plants in cultivation and also among the easiest to grow. While their wild ancestor has not been identified and their origin is unknown, charred legumes of a possible wild-type progenitor have been identified at the Natufian site of the el-Wad Terrace. Carbonised domestic faba bean remains were discovered at three adjacent Neolithic sites in Israel's Lower Galilee (Yiftah'el, Ahi'hud and Nahal Zippori). Based on the radiocarbon dating of these remains, scientists now believe that the domestication of the crop may have begun as early as 8,250 BCE.
Broad beans are still often grown as a cover crop to prevent erosion because they can overwinter and, as a legume, they fix nitrogen in the soil. The broad bean has high plant hardiness; it can withstand harsh and cold climates. Unlike most legumes, the broad bean can be grown in soils with high salinity, as well as in clay soil. However, it prefers rich loams.
In much of the English-speaking world, the name "broad bean" is used for the large-seeded cultivars grown for human food, while "horse bean" and "field bean" refer to cultivars with smaller, harder seeds that are more like the wild species and used for animal feed, though their stronger flavour is preferred in some human food recipes, such as falafel. The name "broad bean" is the most common name in Commonwealth countries like the UK, Canada, Australia and New Zealand, while the term "fava bean" (from for the bean) is used in the United States.
Pests and diseases
Many diseases appear at a higher rate in higher humidity. Therefore, cultivars being bred for higher density should be evaluated for disease problems. This can be mitigated by west–east rows for more sun drying effect.
Disease tolerance is an important part of breeding V. faba.
If transplanted instead of direct seeded there is a lower risk of some diseases including Botrytis fabae.
Parasites
In mainland Europe and North Africa, the plant parasite Orobanche crenata (carnation-scented broomrape) can cause severe impacts on fields of broad beans, devastating their yields.
Fungal diseases
Botrytis fabae
Vicia faba is attacked by Botrytis fabae, the chocolate spot fungus, which can have a severe impact on yield. It is one of the worst diseases in
broad beans, as it results in foliar damage, reduced photosynthesis, and reduced bean productivity. The fungus switches from non-aggressive growth to aggressive pathogenicity under the combination of increased temperature and humidity, which is worsened by low soil potassium and phosphorus content and by the higher humidity caused by higher seeding rates. The non-aggressive phase is marked by small red-brown leaf lesions, and sometimes the same on stems and pods. Treatment is less effective than prevention. Early planting avoids the problematic combination of high temperature and humidity in late spring into early summer. Decreasing seeding rate or thinning after emergence is also effective. Foliar fungicide is effective. If broad beans flower during the height of summer temperatures there is an increased risk of this disease. If transplanted instead of direct seeded there is a lower risk of Botrytis fabae outbreaks.
Erysiphe cichoracearum
Erysiphe cichoracearum overwinters on residue and has alternate hosts. Resistant cultivars and overhead irrigation are preventative. Sulfur fungicides are recommended in severe outbreak.
Fusarium solani
This soil borne pathogen is mitigated by lower temperature, aeration, drainage, and sufficient nutrition. Symptoms include stunting, yellowing, necrotic basal leaves, and brown or red or black streak-shaped root lesions that grow together and may show above the soil as the disease progresses.
Uromyces viciae-fabae var. viciae-fabae
Faba bean rust is a fungal pathogen commonly affecting broad bean plants at maturity, causing small orange dots with yellow halos on the leaves, which may merge to form an orange lawn on both leaf surfaces.
Sclerotinia stem rot
Both Sclerotinia sclerotiorum and S. trifoliorum are pathogens of interest. Lithourgidis et al. have done extensive work over the years, including in 2007 for S. t., 2005 for S. s., and 1989 regarding procedures for field testing with S. s.
Bacterial diseases
Xanthomonas campestris and X. axonopodis
Xanthomonas campestris and X. axonopodis can be inoculated by seed contamination and by overwintering in crop residue. Increased incidence with higher temperatures, rainfall, and humidity. Produces deliquescent, necrotic lesions, sometimes with a wider yellow lesion around them, and in advanced disease the plant will look burned. Can be prevented or treated by use of uninfected seed, resistant cultivars, seed treatments, and copper bactericides.
Pseudomonas syringae
Pseudomonas syringae overwinters on residue. Uninfected seed, rotation, and removal of residue are preventative.
Viral diseases
Faba bean necrotic yellows virus which it shares with other Vicia. Timchenko et al. 2006 find Clink is not obviously necessary but highly conserved nonetheless, suggesting it is maintained by necessity for infection of other Vicia.
Insect pests
Aphis fabae
Broad bean plants are highly susceptible to early summer infestations of the black bean aphid, which can cover large sections of growing plants with infestations, typically starting at the tip of the plant. Severe infestations can significantly reduce yields, and can also cause discolouration of pods and reduction in their saleable values.
Aphis fabae is a major pest. May infest transplants. Reflective plastic mulch may be preventative. May be mechanically removed by high pressure water once plant is established. V. fabae is tolerant to low and medium degrees of infestation, so insecticide application is only required under high infestation.
Toxicity
Beans generally contain phytohaemagglutinin, a lectin that occurs naturally in plants, animals, and humans. Most of the relatively low toxin concentrations found in V. faba can be destroyed by boiling the beans for 10 minutes.
Broad beans are rich in levodopa, and should thus be avoided by those taking irreversible monoamine oxidase inhibitors to prevent a pressor response.
Genetic predisposition
Sufferers of favism must avoid broad beans, as they contain the alkaloid glycoside vicine which may initiate a hemolytic crisis. A low-content vicine-convicine faba bean line was identified in the 1980s and the trait has been introduced into several modern cultivars. Low vicine-convicine faba beans are safe for consumption by G6PD-deficient individuals. As of 2019, a molecular marker may be used for marker-assisted breeding to reduce levels of vicine-convicine in broad beans.
Uses
Culinary
Raw mature broad beans are 11% water, 58% carbohydrates, 26% protein, and 2% fat. A 100-gram reference amount supplies of food energy and numerous essential nutrients in high content (20% or more of the Daily Value, DV). Folate (26% DV), and dietary minerals, such as manganese, phosphorus, magnesium, and iron (range of 52 to 77% DV), have considerable content. B vitamins have moderate to rich content (19 to 48% DV). Broad beans present the highest protein-to-carbohydrate ratio among other popular pulse crops, such as chickpea, pea and lentil. Moreover, their consumption is recommended along with cereals as both foods are complementary in supplying all essential amino acids.
Broad beans are generally eaten while still young and tender, enabling harvesting to begin as early as the middle of spring for plants started under glass or overwintered in a protected location, but even the main crop sown in early spring will be ready from mid to late summer. Horse beans, left to mature fully, are usually harvested in the late autumn, and are then eaten as a pulse. The immature pods are also cooked and eaten, and the young leaves of the plant can also be eaten, either raw or cooked as a pot herb (like spinach).
Preparing broad beans involves first removing the beans from their pods, then steaming or boiling the beans, either whole or after parboiling them to loosen their exterior coating, which is then removed. The beans can be fried, causing the skin to split open, and then salted and/or spiced to produce a savory, crunchy snack.
Algeria
In south Algerian cuisine, broad beans are used to make besarah and doubara. Doubara is popular in the city of Biskra.
China
In the Sichuan cuisine of China, broad beans are combined with soybeans and chili peppers to produce a spicy fermented bean paste called doubanjiang.
Colombia
Broad beans (Colombia: Haba(s)) are a common food in most regions of Colombia, mostly in Bogota and Boyacá.
Ecuador
Steamed broad beans (known as habitas) with cheese is common in the cold-weather regions of Ecuador, especially around the Andes mountains and surroundings of Ambato.
Egypt
Broad beans (Egyptian Masri: ) are a common staple food in the Egyptian diet, eaten by rich and poor alike. Egyptians eat broad beans in various ways: they may be shelled and then dried, or bought dried and then cooked in water on very low heat for several hours. They are the primary ingredient in Egyptian-style falafel (unlike the Levantine style, where the primary ingredient is chickpeas). The most popular way of preparing them in Egypt is by taking the cooked and partially mashed beans and adding oil, salt, and cumin to them. The dish, known as ful medames, is traditionally eaten with bread (generally at breakfast) and is considered one of Egypt's national dishes.
Ethiopia
Broad beans () are one of the most popular legumes in Ethiopia. They are tightly coupled with every aspect of Ethiopian life. They are mainly used as an alternative to peas to prepare a flour called shiro, which is used to make shiro wot (a stew used widely in Ethiopian dishes). During the fasting period in the Ethiopian Orthodox Church tradition called Tsome Filliseta, Tsome arbeå, Tsome Tahsas, and Tsome Hawaria (which are in August, end of February, April, mid-November, beginning of January, and June–July), two uncooked spicy vegetable dishes are made using broad beans. The first is hilibet, a thin, white paste of broad bean flour mixed with pieces of onion, green pepper, garlic, and other spices. The second is siljo, a fermented, sour, spicy thin yellow paste of broad bean flour. Both are served with other stews and injera (a pancake-like bread) during lunch and dinner.
Baqella nifro (boiled broad beans) are eaten as a snack during some holidays and during a time of mourning. This tradition goes well into religious holidays, too. On the Thursday before Good Friday (in the Ethiopian Orthodox Church tradition, tselote hamus (the Prayer of Thursday)), people eat a different kind of nifro called gulban. Gulban is made of peeled half-beans collected and well-cooked with other grains such as wheat, peas, and chickpeas.
England
In England, broad beans are usually boiled.
There is a project aimed at increasing broad bean consumption, particularly by use of broad bean flour in bread.
Finland
In Finnish, the word for "broad bean" is (literally "ox bean"). Broad beans are used to make a meat substitute called Härkis.
Greece
The Greek word fáva (φάβα) does not refer to broad beans, but to the yellow split pea and also to another legume, Lathyrus clymenum. Broad beans are known instead as koukiá (), and are eaten in a stew combined with artichokes, while they are still fresh in their pods. Dried broad beans are eaten boiled, sometimes combined with garlic sauce (skordalia).
In Crete, fresh broad beans are shelled and eaten as a companion to tsikoudia, the local alcoholic drink.
Favism is quite common in Greece because of malaria endemicity in previous centuries, and people afflicted by it do not eat broad beans.
India
In India, broad beans are eaten in the Northeastern state of Manipur. They are locally known as hawai-amubi and are ingredients in the dish eromba.
Iran
Broad beans, or "Baghalee" () are primarily cultivated in the central and north parts of Iran. The city of Kashan has the highest production of broad beans with high quality in terms of the taste, cooking periods and colour. However, broad beans have a very short season (roughly two weeks). The season is usually in the middle of spring. When people have access to fresh beans in season, they cook them in brine and then add vinegar and Heracleum persicum depending on taste. They also make an extra amount to dry to be used year-round. The dried beans can be cooked with rice, which forms one of the most famous dishes in north of Iran (Gilan) called baghalee polo () which means "rice with broad beans". In Iran, broad beans are cooked, served with Golpar-origan and salt and sold on streets in the winter. This food is also available preserved in metal cans.
Iraq
Broad beans which are called Bagilla (باگله/باقله) in the Iraqi dialect of Arabic are a common ingredient in many Iraqi foods. One of the most popular Iraqi dishes that uses the broad bean is Bagilla Bil-Dihin () also called Tishreeb Bagilla (). This dish is a common breakfast dish in Iraq and consists of bread soaked in boiled broad beans’ water then topped with broad beans, melted Ghee, and often also a boiled or fried egg. Fool () is another common breakfast dish in Iraq as well as many other Arab countries and consists of mashed broad beans. Another famous Iraqi dish is Timmen Bagilla (), which is Arabic for 'broad bean rice'. This classic Iraqi dish consists of rice cooked with broad bean and dill.
Italy
In Sardinia, broad beans are traditionally cooked with lard, often substituted or paired with bacon or minced pork.
In Rome, broad beans are popular either cooked with guanciale or with globe artichokes, as side dish together with lamb or kid, or raw with pecorino romano. Fave e pecorino is the traditional dish for 1 May picnic in Liguria, Tuscany, Marche, Umbria and Latium.
In Sicily, maccu is a Sicilian soup prepared with broad beans as a primary ingredient.
In Apulia, broad bean purée with wild chicory is a typical dish.
Japan
Broad beans, called Soramame (Japanese: ) lit: "Sky Bean", are consumed in a variety of ways in Japan. Most commonly, the beans are boiled and are eaten straight or added to rice. It is also consumed as a popular snack called "ikarimame" (Japanese: ) lit: "Anchor Bean", where the beans are roasted or fried.
Luxembourg
Judd mat Gaardebounen, or smoked collar of pork with broad beans, is the national dish of Luxembourg.
Malta
They are a primary ingredient of the Maltese kusksu, a vegetable soup primarily containing broad beans and pasta beads. They are also used in an appetizer called bigilla where they are served as a pureé mixed with olive oil, lemon juice, garlic, parsley and mint. It is served with bread or crackers.
Mexico
In Mexico, broad beans are often eaten in a soup called , meaning "broad bean soup". They are also eaten fried, salted, and dried, as a snack, either by themselves or in combination with other salted, dried beans and nuts.
Morocco
In Morocco, broad beans are cooked, steamed or made into tabiṣart, a dip sold as a street food and commonly eaten in winter.
Nepal
In Nepal, broad beans are called bakulla. They are eaten as a green vegetable when the pods are young, generally stir-fried with garlic. When dried, broad beans are eaten roasted, or mixed with other legumes, such as moong beans, chick peas, and peas, and called qwati. The mixture, soaked and germinated, is cooked as soup and consumed with rice or beaten rice on the occasion of Janai Purnima also known as Rakshya Bandhan, a festival celebrated by the Hindus. The dry and stir-fried version of qwati is called biraula. The qwati soup is believed to reinvigorate the body affected by monsoon paddy season.
Netherlands
In the Netherlands, they are traditionally eaten with fresh savory and some melted butter. The combination of the beans tossed with crispy fried bacon is also common. When rubbed, the velvet insides of the pods are a folk remedy against warts.
Peru
Broad beans (Peruvian Spanish: haba(s)) are eaten fresh or dried toasted, boiled, roasted, stewed or in soup. Habas are one of the essential ingredients of "Pachamanca" in the Andes of Peru, and are also an additive for "Panetela", which is a homemade remedy to keep your child fed and hydrated in cases of diarrhea or stomach infection and even for cholera treatment.
Peruvian dishes with broad beans include:
Aji de habas
Saltado de habas
El chupe de habas
Ajiaco de Papas y habas
Pachamanca
Guiso de habas
Shambar (heavy soup, traditional in Trujillo)
Portugal
Broad beans () are widely cultivated in Portugal and are very popular throughout the country. The most popular dish cooked with favas is "favada", a stew with onion and pork—depending on the region of the country the pork may be chorizo, bacon, pork shoulder, ribs or the mixture of many of these. In Alentejo a lot of coriander will be added in the end. Besides favada, broad beans may be served dry and fried as an appetiser.
Serbia
Broad bean aspic (Serbian: ) is a Serbian winter dish in which the pureed cooked beans are combined with crushed garlic and set in a mould, topped with ground paprika in hot oil.
Spain
Broad beans () are widely cultivated in Spain. Culinary uses vary among regions, but they can be used as the main pulse in a stew (Habas estofadas, michirones) or as an addition to other dishes (menestra, paella). In certain regions they can be eaten while unripe or fried and packaged as a snack.
Sudan
Broad beans are one of the most widely consumed foods in Sudan. For most Sudanese they form the main dish during breakfast time (fatoor), especially more so for city and urban dwellers. The beans are cooked by steadily boiling over a sustained period of time. Similar to Egypt, the cooked beans are mashed, and prepared by adding salt and pepper. For additional flavour, sesame oil is added along with a sprinkling of jibna ("feta" cheese) on top. The dish is then eaten with bread, sometimes mix all in one dish this called (fatta or boash).
Sweden
Broad beans (), which in Sweden were traditionally eaten as soaked brown, and boiled, dried broad beans fried in lard, were for a very long time popular to add to other foods as a filling side, specially with fried pork. The green, raw, and lightly boiled broad beans were used seasonally as a side green.
Syria
In Syria, broad beans are prepared in multiple ways for breakfast, lunch or dinner. Ful medames is the same as the Egyptian dish (it is not mashed though) but with the addition of tomato, parsley and onion and with olive oil. Another version of it includes the addition of tahini (sesame paste), olive oil, garlic and lemon. For lunch, broad beans are cooked with a mix of minced and big chunks of meat and is topped on white rice and eaten with cold yogurt and cucumber salad. Bulgur is sometimes used in preparing this recipe instead of rice. Broad beans are cooked with pieces of garlic, meat and meat stock with the addition of lemon juice and cilantro. This dish is called foulieh and is eaten on the side with rice. The same recipe is prepared without meat as a vegan dish eaten on Lent by Christians in Syria.
Turkey
In Turkey, broad beans are called . This is also the name of a zeytinyağlı dish made by simmering young and tender broad bean pods with chopped onions in olive oil. It is traditionally garnished with dill and served cool, together with yoghurt. Another dish is , a meze prepared by pureeing beans with olive oil. Broad beans are also cooked with artichoke (), which is another zeytinyağlı dish.
Vietnam
In Southern Vietnam, broad beans () are usually stir fried with rice noodles, durians, shrimps, Thai basil, quail eggs and pig intestines in a dry stew called hủ tiếu lòng heo.
Other uses
In ancient Greece and Rome, beans were used in voting; a white bean was used to cast a yes vote, and a black bean for no. Even today, the word koukia (κουκιά) is used unofficially, referring to the votes. Beans were used as a food for the dead, such as during the annual Lemuria festival.
The ancient Roman family name Fabius and the modern political term Fabian derive from this particular bean.
Both Porphyry and Iamblichus report that Pythagoras once persuaded a bull not to eat beans
In Ubykh culture, throwing beans on the ground and interpreting the pattern in which they fall was a common method of divination (favomancy), and the word for "bean-thrower" in that language has become a generic term for seers and soothsayers in general.
The colloquial expression 'not worth a hill of beans' alludes to their widespread economy and association with the peasant diet.
In Italy, broad beans are traditionally sown on 2 November, All Souls Day. Small cakes made in the shape of broad beans (though not out of them) are known as fave dei morti or "beans of the dead". According to tradition, Sicily once experienced a failure of all crops other than the beans; the beans kept the population from starvation, and thanks were given to Saint Joseph. Broad beans subsequently became traditional on Saint Joseph's Day altars in many Italian communities. Some people carry a broad bean for good luck; some believe that if one carries a broad bean, one will never be without the essentials of life. In Rome, on the first of May, Roman families traditionally eat fresh broad beans with Pecorino Romano cheese during a daily excursion in the Campagna. In northern Italy, on the contrary, broad beans are traditionally fed to animals—and so some people, especially the elderly, might frown on human consumption. But in Liguria, a maritime region near northern Italy, broad beans are loved raw, and consumed fresh in early spring as the first product of the garden, alone or with fresh Pecorino Sardo or with local salami from Sant'Olcese. In some Central Italian regions, a once-popular and recently rediscovered fancy food is the bagiana, a soup of fresh or dried broad beans seasoned with onions and beet leaves stir-fried, before being added to the soup, in olive oil and lard (or bacon or cured ham fat).
In Portugal and Spain a Christmas cake called bolo Rei in Portuguese and roscón de reyes in Spanish (King's cake) is baked with a broad bean inside. Whoever eats the slice containing it, is supposed to buy next year's cake.
A similar tradition exists in France, where the fève (originally a dried bean, but often now a small china or metal trinket) is placed in the galette des rois; the person who finds it in their slice becomes the king or queen of the meal, and is often expected to serve the other guests to drink.
Pliny claimed they acted as a laxative.
European folklore also claims that planting beans on Good Friday or during the night brings good luck.
Frederick E Rose (London) Ltd v William H Pim Junior & Co Ltd [1953] 2 QB 450, is an English contract law case where the two litigants had both mistaken feveroles for ordinary horse beans.
Can be used as a green manure, due to nitrogen fixation it produces.
In the Netherlands, roasted or fried broad beans are regarded as a local delicacy of the city of Groningen, and is locally called molleboon. Until the 1800s, the city council used mollebonen for the voting process, sometimes real beans, sometimes made of stone or clay. The word Molleboon became a nickname for the inhabitants of the city.
Research
The first experimental demonstration that the pattern of replication of eukaryotic chromosomes follows the semiconservative DNA replication scheme proposed in 1953 by Watson and Crick was reported in 1957 using V. faba root cells.
Gallery
| Biology and health sciences | Fabales | null |
264502 | https://en.wikipedia.org/wiki/Shed | Shed | A shed is typically a simple, single-storey roofed structure, often used for storage, for hobbies, or as a workshop, and typically serving as outbuilding, such as in a back garden or on an allotment. Sheds vary considerably in their size and complexity of construction, from simple open-sided ones designed to cover bicycles or garden items to large wood-framed structures with shingled roofs, windows, and electrical outlets. Sheds used on farms or in the industry can be large structures. The main types of shed construction are metal sheathing over a metal frame, plastic sheathing and frame, all-wood construction (the roof may be asphalt shingled or sheathed in tin), and vinyl-sided sheds built over a wooden frame. Small sheds may include a wooden or plastic floor, while more permanent ones may be built on a concrete pad or foundation. Sheds may be lockable to deter theft or entry by children, domestic animals, wildlife, etc.
Etymology
The word is recorded in English since 1481, as , possibly a variant of shade. The word shade comes from the Old English word sceadu, which means "shade, shadow, darkness". The term's P.Gmc. cognate, skadwo also means "shady place, protection from glare or heat".
The Old English word is spelled in different ways, such as , shad or shedde, all of which come from an Old Teutonic/Anglo-Saxon root word for separation or division. The first attested usage of the word, in 1481, was in the sentence, .
The Anglo Saxon word shud, which means "cover" may also have been part of the development of the word. In 1440, a shud was defined as a .
Terminology
Depending on the region and type of use, a shed may also be called a shack, outhouse, or "outbuilding". Sheds may be classified as "accessory buildings" in municipal bylaws which may regulate their size, appearance, and distance from the principal building and boundary lines.
Uses
Agricultural sheds
Arena sheds may have a simple open roof structure, or be partially walled or fully enclosed. They are typically used as horse-riding equestrian venues, providing all-year usage of the facility with protection from the weather.
Farm sheds and other outbuildings are used to store farm equipment, tractors, tools, hay, and supplies, or to house horses, cattle, poultry or other farm animals. Run-in sheds are three-sided structures with an open face used for horses and cattle.
Shearing sheds can be large sheds found on sheep stations to accommodate large-scale sheep shearing.
Bike sheds
Bicycle sheds usually contain a bicycle parking rack on which bikes can be supported and locked and a roof to keep rain and/or snow off the bikes and their riders while mounting and dismounting. Bike sheds range from little more than a supported roof to more complex structures with walls and locking doors or gates.
Boatsheds
Boatsheds (or boat sheds) are typically lockable wooden sheds built near a body of water to store small private boats, bathing suits, life vests and related items. Boat sheds used for rowing clubs are generally larger structures for storing rowing skiffs.
Garden sheds
Garden tool sheds, including allotment sheds, are used to store seeds, soil, hoses, portable sprinklers, or garden tools such as hand rakes, shovels, lawnmowers, etc.
Railway sheds
Engine sheds are structures used for the maintenance or storage of railway locomotives. In Britain, these are also called motive power depots.
Goods sheds are railway buildings designed for storing goods before or after carriage in a train.
Train sheds are buildings adjacent to a railway station where the tracks and platforms are covered by a roof. The first train shed was built in 1830 at Liverpool's Crown Street Station.
Snow sheds are strongly built timber or reinforced-concrete tunnels that protect railroad tracks (or roads) from avalanches.
Storage sheds
These may contain any items any person wishes to store and to organise and/or protect from the weather and theft.
Tool sheds
These may contain hand tools and/or power tools used to repair automobiles or for construction.
Wood sheds
These sheds are used for the storage of large quantities of firewood. Woodsheds help protect firewood from adverse weather and moisture, especially in snowy or wet climates. Woodsheds are commonly in close proximity to buildings heated by a wood-burning stove, such as a log cabin.
In the United States, "the woodshed" was the traditional location for parents to administer corporal punishment to children.
Miscellaneous
In the 19th century military barracks, sheds were used as auxiliary buildings for various purposes. The Royal Artillery park barracks in Halifax used sheds as gun sheds, carriage sheds, repair sheds, wheel sheds, wagon sheds and storage sheds.
Construction
Small domestic
The simplest and least-expensive sheds are available in kit form. These kits are designed for regular people to be able to assemble themselves using commonly available tools (e.g., screwdriver). Both shed kits and DIY (do-it-yourself) plans are available for wooden and plastic sheds. Sheds are used to store home and garden tools and equipment such as lawn tractors, and gardening supplies. In addition, sheds can be used to store items that are not suitable for indoor storage, such as petrol (gasoline), pesticides, or herbicides. For homes with small gardens or modest storage needs, there are several types of very small sheds. The sheds not only use less ground area but also have a low profile less likely to obstruct the view or clash with the landscaping.
These small sheds include corner sheds, which fit into a corner (—tall, wide, deep), vertical sheds (), horizontal sheds (), and tool sheds. When a shed is used for tool storage, shelves and hooks are often used to maximise the storage space. Gambrel-style roofed sheds (sometimes called baby barns), which resemble a Dutch-style barn, have a high sloping roofline which increases storage space in the "loft" area. Some Gambrel-styles have no loft and offer the advantage of reduced overall height. Another style of small shed is the saltbox-style shed.
Many sheds have either a pent or apex roof shape. A pent shed features a single roof section that is angled downwards to let rainwater run off, with more headroom at the front than the back. This is a simple, practical design that will fit particularly well next to a wall or fence. It is also usually lower than the typical apex shed, so could be a better choice if there are any height restrictions. A pent shed may be free-standing or attached to a wall (when it is known, unsurprisingly, as a wall shed).
An apex shed has a pointed roof in an inverted V shape similar to the roofline of many houses. Two roof sections meet at a ridge in the middle, providing more headroom in the centre than at the sides. This type is generally regarded as a more attractive and traditional design and may be preferable if the shed is going to be visible from the house.
A twist on the standard apex shape is the reverse apex shed. In this design, the door is set in a side wall instead of the front. The main advantage of the reverse apex design is that the door opens into the widest part of the shed instead of the narrowest, so it is easier to reach into all areas to retrieve or store equipment.
Larger domestic
Larger, more-expensive sheds are typically constructed of wood and include features typically found in house construction, such as windows, a shingled roof, and electrical outlets. Larger sheds provide more space for engaging in hobbies such as gardening, small engine repair, or tinkering. Some sheds have small porches or include furniture, which allows them to be used for relaxation purposes. In some cases, remote workers who live in mild climates use small to medium-sized wooden garden sheds as outdoor offices. There is a growing industry in providing "off the peg" garden offices to cater to this demand, particularly in the UK but also in the US.
Shed owners can customise wooden sheds to match the features (e.g., siding, trim, etc.) of the main house. A number of decorative options can be added to sheds, such as dormers, shutters, flower boxes, finials, and weathervanes. As well, practical options can be added such as benches, ramps, ventilation systems (e.g., in cases where a swimming pool heater is installed in a shed), and electric lighting. Sheds designed for gardening, called "potting sheds", often feature windows or skylights for illumination, ventilation grilles, and a potter's bench for mixing soil and re-potting plants.
Materials
The main types of shed construction are metal sheathing over a metal frame, plastic sheathing and frame, all-wood construction (wood frame, wood siding and wood roof), and vinyl-sided sheds built over a wooden frame. Each type has various advantages and disadvantages that a homeowner has to consider. For example, while metal sheds are fire and termite-resistant, they can rust over time, or be severely damaged by high winds or heavy snow loads. Wood sheds are easier to modify or customise than plastic or metal because carpentry tools and basic carpentry skills are more readily available. Vinyl-sided, wood-framed sheds blend the strength of a wood frame with the maintenance-free aspect of vinyl siding (it does not need to be painted or varnished). The International Building Code (IBC) defines a shed as a building or structure of an accessory character; it classifies them under utility and miscellaneous group U (Chapter 3 Section 312).
Metal
Metal sheds are made from thin sheet metal sheathing (galvanised steel, aluminium, or corrugated iron) attached to a metal frame. Metal sheds are a good choice when long-term strength and resistance to fire, rot, or termites are desired. However, metal sheds may rust over time, particularly if they are constructed from steel that is not galvanised. Be aware that concrete is highly corrosive so care needs to be taken when assembling your shed to avoid contact with the outside panels.
As well, some types of metal sheds that have thin walls are easily dented, which may make some types of thin metal sheds a poor choice for vandal-prone areas or for high-traffic activities such as small businesses. In cold climates, metal sheds with thin walls need to have snow and ice cleared from the roof, because the thin metal may be damaged by a heavy accumulation. Since thin metal sheds weigh much less than wood or PVC plastic sheds, thin metal sheds are more at risk of being damaged by heavy winds. To prevent wind damage, thin metal sheds should be attached to a concrete foundation with screws. In countries where the climate is generally mild, such as Australia, very large metal sheds are used for many types of industry. Corrugated metal sheds may be better able to withstand wind and snow loads, as the corrugated shape makes the metal stronger than flat tin.
Plastic
Plastic shed kits utilising heavy moulded plastics such as PVC and polyethylene may be less expensive than sheet-metal sheds. PVC resins and high-impact, UV light-resistant polyethylene make plastic outdoor sheds stronger, lighter, more durable, and more resistant to denting and chipping than wood, and tend to be more stable. Plastic shed kits sided with vinyl are typically among the least-expensive types of shed construction. Higher-quality sheds use UV-resistant plastic and powder-coated metal frames. Many plastic sheds are modular to allow for easy extensions, peg-boards, shelving, attic-storage, windows, skylights, and other accessories to be added later if these additions are purchased from the manufacturer.
Plastic sheds are not susceptible to termite or wood-boring insect damage, and they require little maintenance. Being rot-proof they do not need to have preservatives applied. This makes them preferable in climates where the weather can be changeable, such as the United Kingdom. Unlike wooden or metal sheds, which often require a permit to build, in many areas, plastic sheds do not. However, this is something property owners will need to verify. A call to your council/town's planning or building code office can provide information on permits.
Wood
Wooden sheds have a natural look that can blend in well with garden environments. Despite the strength of wood, over time, untreated and neglected wood can rot, split, warp or become susceptible to mold and mildew, so wood sheds should be treated for protection with stain and varnish. Wood sheds need regular maintenance. This includes keeping plant matter and debris from piling up beside the walls and on the roof, and occasional rot-proofing with preservatives. Sheds are sometimes also re-stained or varnished at times for aesthetic and wood protection reasons. Fire and, in some regions, termite attack are also potential problems.
Stains and preservatives can be applied to wood sheds to prevent damage to the wood caused by exposure to rain, damp ground, UV light, harsh climatic conditions, fungal attack and wood-boring insects. If a coloured preservative oil or stain is used, a wooden shed can either be made to stand out as a feature within a garden, or to blend in with its surroundings. Red cedar coloured stain is popular. Legislation such as the European Biocidal Products Regulation has reduced the number of effective active ingredients available for wood preservative formulations. For this reason, in recent years, there has been a greater emphasis on preserving wood by keeping it dry, for example through the application of water-repellent "wood protection creams."
Some types of wood, such as cedar, are more naturally resistant to water damage.
When looking for a wooden shed, it is important to understand the difference between the two types of preservative used in their manufacture. The timber will have been treated in one of two ways: dip treatment or pressure treatment. Dip-treated sheds are made from components that are lowered into a tank of preservatives before the panels are assembled. This is a quick and simple process that keeps costs down and encourages manufacturers to produce a wide variety, making dip-treated sheds the most popular and affordable type on the market. They are easily recognisable by their golden brown colour, which is due to a dye added to the preservative. Most manufacturers offer a 10-year anti-rot guarantee on dip-treated sheds, but they have to be re-coated every year or two.
Pressure-treated sheds are made from timber planks that have had the moisture sucked out of them under vacuum conditions in a special cylinder. A powerful preservative is then forced into the wood at high pressure until it is absorbed deep into the grain, becoming an integral part of the timber. This provides excellent protection against the weather—so much so that manufacturers generally give a 15-year anti-rot guarantee. These sheds are usually distinguished by a pale green tinge which will fade eventually to a silvery grey. Although pressure-treated sheds tend to be more expensive than dip-treated ones, their big advantage is that they will not need any further preservative treatment during the guarantee period, saving owners time and money.
One advantage of using wood sheds over metal versions is that it is easier to modify them by adding windows, doors, shelving, or exterior trim (etc.) because wood can be cut and drilled using commonly available tools, whereas a plastic or metal shed requires specialised tools. Some homeowners may prefer wood sheds because wood is a renewable resource.
Vinyl-sided
Vinyl-sided sheds are typically built with standard wood framing construction and oriented strand board (OSB) on the walls covered with standard vinyl siding. The vinyl siding protects the OSB wood and the frame from moisture from rain and snow. Vinyl-sided sheds never need to be painted, and are maintenance-free. They are stronger than plastic or metal sheds, and are usually built to conform with the local building codes. They offer good value for money because they hold up in all weather, including winters with heavy snowfall, as they use a strong wooden frame and the OSB panels have stronger structural support than thin metal or PVC siding or roofs. Metal, plastic and resin sheds are cheaper, but they cannot handle the weight of snow in winter (roofs may cave in). Vinyl sheds also offer more colour options.
Asbestos
In the early and middle years of the 20th century, many garden sheds and domestic garages were made of asbestos-cement sheets supported on a very light angle-iron frame. Concerns about safety led to the practice being discontinued, but they were cheap and long-lasting, and many can still be seen in British gardens. Advice on continued use or disposal is available.
Culture
In Australia and New Zealand the term shed can be used to refer to any building that is not a residence and which may be open at the ends or sides, or both. Australia's passion for sheds is documented in Mark Thomson's Blokes and Sheds (1998). Jim Hopkins' similarly titled Blokes & Sheds (1998), with photographer Julie Riley Hopkins, profiles amateur inventors from across New Zealand. Hopkins and Riley followed up that book with Inventions from the Shed (1999) and a 5-part film documentary series with the same name. Gordon Thorburn also examined the shed proclivity in his book Men and Sheds (2002), as did Gareth Jones in Shed Men (2004).
Recently, "Men's Sheds" have become common in Australia. In New Zealand, the bi-monthly magazine The Shed appeals to the culture of "blokes" who do woodwork or metalwork DIY projects in their sheds. The Australian Men's Shed Association is one organisation that has been set up involving sheds.
Another magazine called The Shed, a bi-monthly PDF magazine produced in the UK, but with a global audience, targets people who work (usually in creative industries) in garden offices, sheds and other shed-like atmospheres. In the UK, people have long enjoyed working in their potting sheds; the slang term "sheddie", refers to a person enamoured of shed-building, testifies to the place of sheds in the UK popular culture. The Usenet newsgroup "uk.rec.sheds" has long championed this subculture.
Since 2007 there has been a UK competition called Shed of the Year. Each year British sheddies enter their shed builds and after a short list is produced (including Pub Sheds, Eco Sheds, Workshops & unexpected categories), a public vote helps to decide the ultimate Shed, it also featured on Channel 4 television as George Clarke's Amazing Spaces: Shed of the Year, for four series (with host George Clarke)
In the United States, Shed Builder Magazine is a bimonthly magazine dedicated to the builders, dealers, and manufacturers within the shed industry. The magazine owns and manages Shed Builder Expo, a yearly, two-day conference for the shed industry.
Author Gordon Thorburn examined the shed proclivity in his book Men and Sheds, which argues that a "place of retreat" is a "male necessity" which provides men with solace, especially during their retirement. In contrast, in the novel Cold Comfort Farm by Stella Gibbons, Aunt Ada Doom saw "something nasty in the woodshed" and retreated to her bed for half a century.
To woodshed, or 'shed, in jazz jargon, is "to shut oneself up, away from the world, and practice long and hard, as in 'going to the woodshed'."
A shed built onto the chassis of an old car, and called Fastest Shed, is legally roadworthy in the UK, and holds the world speed record for sheds.
| Technology | Buildings and infrastructure | null |
264583 | https://en.wikipedia.org/wiki/Phosphine | Phosphine | Phosphine (IUPAC name: phosphane) is a colorless, flammable, highly toxic compound with the chemical formula PH3, classed as a pnictogen hydride. Pure phosphine is odorless, but technical grade samples have a highly unpleasant odor like rotting fish, due to the presence of substituted phosphine and diphosphane (). With traces of present, is spontaneously flammable in air (pyrophoric), burning with a luminous flame. Phosphine is a highly toxic respiratory poison, and is immediately dangerous to life or health at 50 ppm. Phosphine has a trigonal pyramidal structure.
Phosphines are compounds that include and the organophosphines, which are derived from by substituting one or more hydrogen atoms with organic groups. They have the general formula . Phosphanes are saturated phosphorus hydrides of the form , such as triphosphane. Phosphine, PH3, is the smallest of the phosphines and the smallest of the phosphanes.
History
Philippe Gengembre (1764–1838), a student of Lavoisier, first obtained phosphine in 1783 by heating white phosphorus in an aqueous solution of potash (potassium carbonate).
Perhaps because of its strong association with elemental phosphorus, phosphine was once regarded as a gaseous form of the element, but Lavoisier (1789) recognised it as a combination of phosphorus with hydrogen and described it as phosphure d'hydrogène (phosphide of hydrogen).
In 1844, Paul Thénard, son of the French chemist Louis Jacques Thénard, used a cold trap to separate diphosphine from phosphine that had been generated from calcium phosphide, thereby demonstrating that is responsible for spontaneous flammability associated with , and also for the characteristic orange/brown color that can form on surfaces, which is a polymerisation product. He considered diphosphine's formula to be , and thus an intermediate between elemental phosphorus, the higher polymers, and phosphine. Calcium phosphide (nominally ) produces more than other phosphides because of the preponderance of P-P bonds in the starting material.
The name "phosphine" was first used for organophosphorus compounds in 1857, being analogous to organic amines (). The gas was named "phosphine" by 1865 (or earlier).
Structure and reactions
is a trigonal pyramidal molecule with C3v molecular symmetry. The length of the P−H bond is 1.42 Å, the H−P−H bond angles are 93.5°. The dipole moment is 0.58 D, which increases with substitution of methyl groups in the series: , 1.10 D; , 1.23 D; , 1.19 D. In contrast, the dipole moments of amines decrease with substitution, starting with ammonia, which has a dipole moment of 1.47 D. The low dipole moment and almost orthogonal bond angles lead to the conclusion that in the P−H bonds are almost entirely and phosphorus 3s orbital contributes little to the P-H bonding. For this reason, the lone pair on phosphorus is predominantly formed by the 3s orbital of phosphorus. The upfield chemical shift of it 31P NMR signal accords with the conclusion that the lone pair electrons occupy the 3s orbital (Fluck, 1973). This electronic structure leads to a lack of nucleophilicity in general and lack of basicity in particular (pKaH = –14), as well as an ability to form only weak hydrogen bonds.
The aqueous solubility of is slight: 0.22 cm3 of gas dissolves in 1 cm3 of water. Phosphine dissolves more readily in non-polar solvents than in water because of the non-polar P−H bonds. It is technically amphoteric in water, but acid and base activity is poor. Proton exchange proceeds via a phosphonium () ion in acidic solutions and via phosphanide () at high pH, with equilibrium constants Kb = and Ka = . Phosphine reacts with water only at high pressure and temperature, producing phosphoric acid and hydrogen:
Burning phosphine in the air produces phosphoric acid:
Preparation and occurrence
Phosphine may be prepared in a variety of ways. Industrially it can be made by the reaction of white phosphorus with sodium or potassium hydroxide, producing potassium or sodium hypophosphite as a by-product.
Alternatively, the acid-catalyzed disproportionation of white phosphorus yields phosphoric acid and phosphine. Both routes have industrial significance; the acid route is the preferred method if further reaction of the phosphine to substituted phosphines is needed. The acid route requires purification and pressurizing.
Laboratory routes
It is prepared in the laboratory by disproportionation of phosphorous acid:
Alternative methods are the hydrolysis zinc phosphide:
Some other metal phosphides could be used including aluminium phosphide, or calcium phosphide. Pure samples of phosphine, free from , may be prepared using the action of potassium hydroxide on phosphonium iodide:
Occurrence
Phosphine is a worldwide constituent of the Earth's atmosphere at very low and highly variable concentrations. It may contribute significantly to the global phosphorus biochemical cycle. The most likely source is reduction of phosphate in decaying organic matter, possibly via partial reductions and disproportionations, since environmental systems do not have known reducing agents of sufficient strength to directly convert phosphate to phosphine.
It is also found in Jupiter's atmosphere.
Possible extraterrestrial biosignature
In 2020 a spectroscopic analysis was reported to show signs of phosphine in the atmosphere of Venus in quantities that could not be explained by known abiotic processes. Later re-analysis of this work showed interpolation errors had been made, and re-analysis of data with the fixed algorithm do not result in the detection of phosphine. The authors of the original study then claimed to detect it with a much lower concentration of 1 ppb.
Applications
Organophosphorus chemistry
Phosphine is a precursor to many organophosphorus compounds. It reacts with formaldehyde in the presence of hydrogen chloride to give tetrakis(hydroxymethyl)phosphonium chloride, which is used in textiles. The hydrophosphination of alkenes is versatile route to a variety of phosphines. For example, in the presence of basic catalysts adds of Michael acceptors. Thus with acrylonitrile, it reacts to give tris(cyanoethyl)phosphine:
Acid catalysis is applicable to hydrophosphination with isobutylene and related analogues:
where R is , alkyl, etc.
Microelectronics
Phosphine is used as a dopant in the semiconductor industry, and a precursor for the deposition of compound semiconductors. Commercially significant products include gallium phosphide and indium phosphide.
Fumigant (pest control)
Phosphine is an attractive fumigant because it is lethal to insects and rodents, but degrades to phosphoric acid, which is non-toxic. As sources of phosphine, for farm use, pellets of aluminium phosphide (AlP), calcium phosphide (), or zinc phosphide () are used. These phosphides release phosphine upon contact with atmospheric water or rodents' stomach acid. These pellets also contain reagents to reduce the potential for ignition or explosion of the released phosphine.
An alternative is the use of phosphine gas itself which requires dilution with either or or even air to bring it below the flammability point. Use of the gas avoids the issues related with the solid residues left by metal phosphide and results in faster, more efficient control of the target pests.
One problem with phosphine fumigants is the increased resistance by insects.
Toxicity and safety
Deaths have resulted from accidental exposure to fumigation materials containing aluminium phosphide or phosphine. It can be absorbed either by inhalation or transdermally. As a respiratory poison, it affects the transport of oxygen or interferes with the utilization of oxygen by various cells in the body. Exposure results in pulmonary edema (the lungs fill with fluid). Phosphine gas is heavier than air so it stays near the floor.
Phosphine appears to be mainly a redox toxin, causing cell damage by inducing oxidative stress and mitochondrial dysfunction. Resistance in insects is caused by a mutation in a mitochondrial metabolic gene.
Phosphine can be absorbed into the body by inhalation. The main target organ of phosphine gas is the respiratory tract. According to the 2009 U.S. National Institute for Occupational Safety and Health (NIOSH) pocket guide, and U.S. Occupational Safety and Health Administration (OSHA) regulation, the 8 hour average respiratory exposure should not exceed 0.3 ppm. NIOSH recommends that the short term respiratory exposure to phosphine gas should not exceed 1 ppm. The Immediately Dangerous to Life or Health level is 50 ppm. Overexposure to phosphine gas causes nausea, vomiting, abdominal pain, diarrhea, thirst, chest tightness, dyspnea (breathing difficulty), muscle pain, chills, stupor or syncope, and pulmonary edema. Phosphine has been reported to have the odor of decaying fish or garlic at concentrations below 0.3 ppm. The smell is normally restricted to laboratory areas or phosphine processing since the smell comes from the way the phosphine is extracted from the environment. However, it may occur elsewhere, such as in industrial waste landfills. Exposure to higher concentrations may cause olfactory fatigue.
Fumigation hazards
Phosphine is used for pest control, but its usage is strictly regulated due to high toxicity. Gas from phosphine has high mortality rate and has caused deaths in Sweden and other countries.
Because the previously popular fumigant methyl bromide has been phased out in some countries under the Montreal Protocol, phosphine is the only widely used, cost-effective, rapidly acting fumigant that does not leave residues on the stored product. Pests with high levels of resistance toward phosphine have become common in Asia, Australia and Brazil. High level resistance is also likely to occur in other regions, but has not been as closely monitored. Genetic variants that contribute to high level resistance to phosphine have been identified in the dihydrolipoamide dehydrogenase gene. Identification of this gene now allows rapid molecular identification of resistant insects.
Explosiveness
Phosphine gas is denser than air and hence may collect in low-lying areas. It can form explosive mixtures with air, and may also self-ignite.
In fiction
Anne McCaffrey's Dragonriders of Pern series features genetically engineered dragons that breathe fire by producing phosphine by extracting it from minerals of their native planet.
In the 2008 pilot of the crime drama television series Breaking Bad, Walter White poisons two rival gangsters by adding red phosphorus to boiling water to produce phosphine gas. However, this reaction in reality would require white phosphorus instead, and for the water to contain sodium hydroxide.
| Physical sciences | Hydrogen compounds | Chemistry |
264589 | https://en.wikipedia.org/wiki/Lapwing | Lapwing | Lapwings (subfamily Vanellinae) are any of various ground-nesting birds (family Charadriidae) akin to plovers and dotterels. They range from in length, and are noted for their slow, irregular wingbeats in flight and a shrill, wailing cry.
The traditional terms "plover", "lapwing", and "dotterel" do not correspond exactly to current taxonomic models; thus, several of the Vanellinae are often called plovers, and one a dotterel, while a few of the "true" plovers (subfamily Charadriinae) are known colloquially as lapwings. In general, a lapwing can be thought of as a larger plover.
In Europe's Anglophone countries, lapwing refers specifically to the northern lapwing, the only member of this group to occur in most of the continent and thus the first bird to go by the English name lapwing (also known as peewit or pyewipe).
In the fanciful taxonomy promoted by medieval courtesy books, a group of lapwings was called a "deceit".
Systematics
While authorities generally agree that there are approximately 25 species of Vanellinae, classifications within the subfamily remain confused. Some workers have gone so far as to group all the "true" lapwings (except the red-kneed dotterel) into the single genus Vanellus. Current consensus favors a more moderate position, but it is unclear which genera to split. The Handbook of Birds of the World provisionally places all Vanellinae in Vanellus except the red-kneed dotterel, which is in the monotypic Erythrogonys. Its plesiomorphic habitus resembles that of plovers, but details like the missing hallux (hind toe) are like those of lapwings: it is still not entirely clear whether it is better considered the most basal plover or lapwing. The IOC also recognizes a monotypic genus Hoploxypterus for the pied plover.
Many coloration details of the red-kneed dotterel also occur here and there among the living members of the main lapwing clade. Its position as the most basal of the living Vanellinae or just immediately outside it thus means that their last common ancestor – or even the last common ancestor of plovers and lapwings – almost certainly was a plover-sized bird with a black crown and breast-band, a white feather patch at the wrist, no hallux, and a lipochromic (probably red) bill with a black tip. Its legs were most likely black or the color of the bill's base.
Evolution
The fossil record of the Vanellinae is scant and mostly recent in origin; no Neogene lapwings seem to be known. On the other hand, it appears as if early in their evolutionary history the plovers, lapwings and dotterels must have been almost one and the same, and they are hard to distinguish osteologically even today. Thus, since the Red-kneed Dotterel is so distinct that it might arguably be considered a monotypic subfamily, reliably dating its divergence from a selection of true lapwings and plovers would also give a good idea of charadriid wader evolution altogether.
A mid-Oligocene – c.28 mya (million years ago) – fossil from Rupelmonde in Belgium has been assigned to Vanellus, but even if the genus were broadly defined, it is entirely unclear if the placement is correct. Its age ties in with the appearance of the first seemingly distinct Charadriinae at about the same time, and with the presence of more basal Charadriidae a few million years earlier. However, the assignment of fragmentary fossils to Charadriinae or Vanellinae is not easy. Thus, it is very likely that the charadriid waders originate around the Eocene-Oligocene boundary – roughly 40–30 mya – but nothing more can be said at present. If the Belgian fossil is not a true lapwing, there are actually no Vanellinae fossils known before the Quaternary.
The Early Oligocene fossil Dolicopterus from Ronzon, France may be such an ancestral member of the Charadriidae or even the Vanellinae, but it has not been studied in recent decades and is in dire need of review.
Apart from the prehistoric Vanellus, the extinct lapwing genus Viator has been described from fossils. Its remains were found in the tar pits of Talara in Peru and it lived in the Late Pleistocene. Little is known of this rather large lapwing; it may actually belong in Vanellus.
The remaining Charadrii are highset and/or chunky birds, even decidedly larger than a lot of the scolopacid waders. The evolutionary trend regarding the Charadriidae – which make up most of the diversity of the Charadrii – thus runs contrary to Cope's Rule.
List of species in taxonomic order
Genus Vanellus
Northern lapwing, Vanellus vanellus
White-headed lapwing, Vanellus albiceps
Southern lapwing, Vanellus chilensis
Grey-headed lapwing, Vanellus cinereus
Crowned lapwing, Vanellus coronatus
Long-toed lapwing, Vanellus crassirostris
River lapwing or spur-winged lapwing, Vanellus duvaucelii
Red-wattled lapwing, Vanellus indicus
Masked lapwing, Vanellus miles
Spur-winged lapwing or spur-winged plover, Vanellus spinosus
Banded lapwing, Vanellus tricolor
Blacksmith lapwing, Vanellus armatus
Black-headed lapwing, Vanellus tectus
Yellow-wattled lapwing, Vanellus malabaricus
Senegal lapwing, Vanellus lugubris
Black-winged lapwing, Vanellus melanopterus
African wattled lapwing, Vanellus senegallus
Spot-breasted lapwing, Vanellus melanocephalus
Brown-chested lapwing, Vanellus superciliosus
Javanese wattled lapwing, Vanellus macropterus
Sociable lapwing, Vanellus gregarius
White-tailed lapwing, Vanellus leucurus
Andean lapwing, Vanellus resplendens
Genus Hoploxypterus
Pied lapwing, Hoploxypterus cayanus
Genus Erythrogonys
Red-kneed dotterel, Erythrogonys cinctus
| Biology and health sciences | Charadriiformes | Animals |
264604 | https://en.wikipedia.org/wiki/Ovenbird%20%28family%29 | Ovenbird (family) | Ovenbirds or furnariids are a large family of small suboscine passerine birds found from Mexico and Central to southern South America. They form the family Furnariidae. This is a large family containing around 321 species and 71 genera. The ovenbird (Seiurus aurocapilla), which breeds in North America, is not a furnariid – rather it is a distantly related bird of the wood warbler family, Parulidae.
The ovenbirds are a diverse group of insectivores which get their name from the elaborate, vaguely "oven-like" clay nests built by the horneros, although most other ovenbirds build stick nests or nest in tunnels or clefts in rock. The Spanish word for "oven" (horno) gives the horneros their name. Furnariid nests are always constructed with a cover, and up to six pale blue, greenish or white eggs are laid. The eggs hatch after 15 to 22 days, and the young fledge after a further 13 to 20 days.
They are small to medium-sized birds, ranging from 9 to 35 cm in length. While individual species often are habitat specialists, species of this family can be found in virtually any Neotropical habitat, ranging from city parks inhabited by rufous horneros, to tropical Amazonian lowlands by many species of foliage-gleaners, to temperate barren Andean highlands inhabited by several species of miners. Two species, the seaside and the surf cinclodes, are associated with rocky coasts.
Taxonomy and systematics
The woodcreepers (formerly Dendrocolaptidae) were merged into this family, following analysis of sequences. While confirming the overall phylogenetic pattern, other scientists instead opted for maintaining the woodcreepers as a separate family, while splitting the ovenbirds (as traditionally defined) into two families, Furnariidae and Scleruridae.
The cladogram below showing the subfamilies of the ovenbirds is based on a molecular genetic studies that revealed that Sclerurinae was the first group to diverge The species numbers are from the list maintained by the International Ornithologists' Union (IOC).
The phylogeny of the Furnariidae is now well understood thanks to multiple analyses of nuclear and mitochondrial DNA. Among other discoveries, the classification of several genera had to be revised. The taxonomic arrangement presented below is based on molecular genetic studies of ovenbird relationships. However, because ovenbirds and woodcreepers are treated here as a single family some taxonomic ranks were modified. For more detail see "List of ovenbird species".
Subfamily: Sclerurinae – miners and leaftossers
Genus Geositta – miners (11 species)
Genus Sclerurus – leaftossers (7 species)
Subfamily: Dendrocolaptinae – woodcreepers
Tribe: Sittasomini – "intermediate" woodcreepers
Genus Dendrocincla – woodcreepers (6 species)
Genus Deconychura – long-tailed woodcreepers (3 species)
Genus Sittasomus – olivaceous woodcreeper
Genus Certhiasomus – spot-throated woodcreeper (genus introduced in 2010 for Deconychura stictolaema)
Tribe: Dendrocolaptini – "strong-billed" woodcreepers
Genus Glyphorynchus – wedge-billed woodcreeper
Genus Nasica – long-billed woodcreeper
Genus Dendrexetastes – cinnamon-throated woodcreeper
Genus Dendrocolaptes – woodcreepers (5 species)
Genus Hylexetastes – woodcreepers (3 species)
Genus Xiphocolaptes – woodcreepers (4 species)
Genus Dendroplex – straight-billed woodcreepers (2 species, formerly in Xiphorhynchus)
Genus Xiphorhynchus – woodcreepers (13 species)
Genus Lepidocolaptes – narrow-billed woodcreepers (11 species)
Genus Drymornis – scimitar-billed woodcreeper
Genus Drymotoxeres – greater scythebill
Genus Campylorhamphus – scythebills (6 species)
Subfamily: Furnariinae – Neotropical ovenbirds and allies
Genus: Xenops – xenops (5 species)
Genus Berlepschia – point-tailed palmcreeper
Tribe Pygarrhichini
Genus Pygarrhichas – white-throated treerunner
Genus Microxenops – rufous-tailed xenops
Genus Ochetorhynchus – earthcreepers (4 species formerly included in Upucerthia)
Tribe Furnariini – horneros and allies
Genus Pseudocolaptes – tuftedcheeks (3 species)
Genus Premnornis – rusty-winged barbtail
Genus Tarphonomus – earthcreepers (genus introduced in 2007 for 2 species formerly included in Upucerthia)
Genus Geocerthia – striated earthcreeper (genus introduced in 2009 for U. serrrana)
Genus Upucerthia – earthcreepers (4 species)
Genus Cinclodes – cinclodes (15 species)
Genus Furnarius – horneros (8 species)
Genus Lochmias – sharp-tailed streamcreeper
Genus Phleocryptes – wren-like rushbird
Genus Limnornis – curve-billed reedhaunter
Tribe Philydorini – foliage-gleaners and allies
Genus Megaxenops – great xenops
Genus Anabazenops – foliage-gleaners (2 species)
Genus Ancistrops – chestnut-winged hookbill
Genus Cichlocolaptes – (2 species)
Genus Heliobletus – sharp-billed treehunter
Genus Neophilydor – foliage-gleaners (genus introduced in 2023 for 2 species formerly in Philydor)
Genus Philydor – foliage-gleaners (3 species)
Genus Dendroma – foliage-gleaners (2 species)
Genus Anabacerthia – foliage-gleaners (5 species)
Genus Syndactyla – foliage-gleaners (8 species)
Genus Clibanornis – (5 species)
Genus Thripadectes – treehunters (7 species)
Genus Automolus – foliage-gleaners (11 species)
Tribe Synallaxini – spinetails and allies
Genus Margarornis – treerunners (4 species)
Genus Premnoplex – typical barbtails (2 species)
Genus Aphrastura – rayaditos (3 species)
Genus Hellmayrea – white-browed spinetail
Genus Sylviorthorhynchus – (2 species)
Genus Leptasthenura – tit-spinetails (9 species)
Genus Phacellodomus – thornbirds (10 species)
Genus Anumbius – firewood-gatherer
Genus Coryphistera – lark-like brushrunner
Genus Pseudoseisura – cacholotes (4 species)
Genus Pseudasthenes – false canasteros
Genus Spartonoica – bay-capped wren-spinetail
Genus Asthenes – canasteros (29 species)
Genus Certhiaxis – spinetails (2 species)
Genus Mazaria – white-bellied spinetail
Genus Schoeniophylax – chotoy spinetail
Genus Synallaxis – spinetails (37 species)
Genus Siptornis – spectacled prickletail
Genus Metopothrix – orange-fronted plushcrown
Genus Xenerpestes – graytails (2 species)
Genus Acrobatornis – pink-legged graveteiro
Genus Limnoctites – reedhaunters (2 species)
Genus Thripophaga – softtails (4 species)
Genus Cranioleuca – typical spinetails (20 species)
Genus Roraimia – Roraiman barbtail
The phylogenetic tree shown below is based on a large-scale genetic 2020 study of the suboscines by Michael Harvey and collaborators. The tawny tit-spinetail (Leptasthenura yanacencis) has been moved to the genus Sylviorthorhynchus, the sulphur-bearded spinetail (Cranioleuca sulphurifera) has been moved to the genus Limnoctites and its English name changed to the sulphur-bearded reedhaunter, and the white-bellied spinetail (Synallaxis propinqua) has been placed in the monotypic genus Mazaria. These changes are included in the tree shown below. The remaining paraphyletic genera are flagged in the tree by an asterisk.
In 2009, the large ovenbird family was divided into tribes by Robert Moyle and collaborators. The tribes as defined in the 2009 article do not fit well with the revised taxonomy of Harvey and are not included here. For example, the tribe Furnariini as defined in the 2009 article is not monophyletic in the Harvey phylogeny. The species numbers in the cladogram are from the list maintained by the International Ornithologists' Union (IOC).
Fossil record
Furnariids boast a notable fossil record for a passerine family. Numerous fossils comprising multiple skeletal elements, including cranial remains, have facilitated the identification and description of five distinct fossil species. Among these, two have been classified within the extant genera Cinclodes and Pseudoseisura, while the remaining three belong into the extinct genus Pseudoseisuropsis. All fossil are of Pleistocene age.
†Pseudoseisuropsis nehuen Noriega 1991, early Pleistocene of Argentina.
†Pseudoseisuropsis cuelloi Claramunt & Rinderknecht 2005, late Pleistocene of Uruguay.
†Pseudoseisuropsis wintu Stefanini et al. 2016, Early Pleistocene of Argentina.
†Cinclodes major Toni 1977, Pleistocene of Argentina.
†Pseudoseisura cursor Toni & Noriega, 2001, Pleistocene of Argentina.
| Biology and health sciences | Tyranni | null |
264714 | https://en.wikipedia.org/wiki/Iowa-class%20battleship | Iowa-class battleship | The Iowa class was a class of six fast battleships ordered by the United States Navy in 1939 and 1940. They were initially intended to intercept fast capital ships such as the Japanese and serve as the "fast wing" of the U.S. battle line. The Iowa class was designed to meet the Second London Naval Treaty's "escalator clause" limit of standard displacement. Beginning in August 1942, four vessels, , , , and , were completed; two more, and , were laid down but canceled in 1945 and 1958, respectively, before completion, and both hulls were scrapped in 1958–1959.
The four Iowa-class ships were the last battleships commissioned in the U.S. Navy. All older U.S. battleships were decommissioned by 1947 and stricken from the Naval Vessel Register (NVR) by 1963. Between the mid-1940s and the early 1990s, the Iowa-class battleships fought in four major U.S. wars. In the Pacific Theater of World War II, they served primarily as fast escorts for s of the Fast Carrier Task Force and also shelled Japanese positions. During the Korean War, the battleships provided naval gunfire support (NGFS) for United Nations forces, and in 1968, New Jersey shelled Viet Cong and Vietnam People's Army forces in the Vietnam War. All four were reactivated and modernized at the direction of the United States Congress in 1981, and armed with missiles during the 1980s, as part of the 600-ship Navy initiative. During Operation Desert Storm in 1991, Missouri and Wisconsin fired missiles and guns at Iraqi targets.
Costly to maintain, the battleships were decommissioned during the post-Cold War drawdown in the early 1990s. All four were initially removed from the Naval Vessel Register, but the United States Congress compelled the Navy to reinstate two of them on the grounds that existing shore bombardment capability would be inadequate for amphibious operations. This resulted in a lengthy debate over whether battleships should have a role in the modern navy. Ultimately, all four ships were stricken from the Naval Vessel Register and released for donation to non-profit organizations. With the transfer of Iowa in 2012, all four are museum ships part of non-profit maritime museums across the US.
Background
The vessels that eventually became the Iowa-class battleships were born from the U.S. Navy's War Plan Orange, a Pacific war plan against Japan. War planners anticipated that the U.S. fleet would engage and advance in the Central Pacific, with a long line of communication and logistics that would be vulnerable to high-speed Japanese cruisers and capital ships. The chief concern was that the U.S. Navy's traditional 21-knot battle line of "Standard-type" battleships would be too slow to force these Japanese task forces into battle, while faster aircraft carriers and their cruiser escorts would be outmatched by the Japanese battlecruisers, which had been upgraded in the 1930s to fast battleships. As a result, the U.S. Navy envisioned a fast detachment of the battle line that could bring the Japanese fleet into battle. Even the new standard battle line speed of 27 knots, as the preceding and battleships were designed for, was not considered enough and during their development processes, designs that could achieve over 30 knots in order to counter the threat of fast "big gun" ships were seriously considered. At the same time, a special strike force consisting of fast battleships operating alongside carriers and destroyers was being envisaged; such a force could operate independently in advance areas and act as scouts. This concept eventually evolved into the Fast Carrier Task Force, though initially the carriers were believed to be subordinate to the battleship.
Another factor was the "escalator clause" of the Second London Naval Treaty, which reverted the gun caliber limit from to . Japan had refused to sign the treaty and in particular refused to accept the 14-inch gun caliber limit or the 5:5:3 ratio of warship tonnage limits for Britain, the United States, and Japan, respectively. This resulted in the three treaty powers, the United States, Britain, and France, invoking the caliber escalator clause after April 1937. Circulation of intelligence evidence in November 1937 of Japanese capital ships violating naval treaties caused the treaty powers to expand the escalator clause in June 1938, which amended the standard displacement limit of battleships from to .
Design
Early studies
Work on what would eventually become the Iowa-class battleship began on the first studies in early 1938, at the direction of Admiral Thomas C. Hart, head of the General Board, following the planned invocation of the "escalator clause" that would permit maximum standard capital ship displacement of . Using the additional over previous designs, the studies included schemes for "slow" battleships that increased armament and protection as well as "fast" battleships capable of or more. One of the "slow" designs was an expanded South Dakota class carrying either twelve 16-inch/45 caliber Mark 6 guns or nine /48 guns and with more armor and a power plant large enough to drive the larger ship through the water at the same 27-knot maximum speed as the South Dakotas. While the "fast" studies would result in the Iowa class, the "slow" design studies would eventually settle on twelve 16-inch guns and evolve into the design for the after all treaty restrictions were removed following the start of World War II. Priority was given to the "fast" design in order to counter and defeat Japan's Kongō-class fast battleships, whose higher speed advantage over existing U.S. battleships might let them "penetrate U.S. cruisers, thereby making it 'open season' on U.S. supply ships", and then overwhelm the Japanese battle line was therefore a major driving force in setting the design criteria for the new ships, as was the restricting width of the Panama Canal.
For "fast" battleships, one such design, pursued by the Design Division section of the Bureau of Construction and Repair, was a "cruiser-killer". Beginning on 17 January 1938, under Captain A.J. Chantry, the group drew up plans for ships with twelve 16-inch and twenty guns, Panamax capability but otherwise unlimited displacement, a top speed of and a range of when traveling at the more economical speed of . Their plan fulfilled these requirements with a ship of standard displacement, but Chantry believed that more could be done if the ship were to be this large; with a displacement greater than that of most battleships, its armor would have protected it only against the weapons carried by heavy cruisers.
Three improved plans – "A", "B", and "C" – were designed at the end of January. An increase in draft, vast additions to the armor, and the substitution of twelve guns in the secondary battery were common among the three designs. "A" was the largest, at standard, and was the only one to still carry the twelve 16-inch guns in four triple turrets (3-gun turrets according to US Navy). It required to make . "B" was the smallest at standard; like "A" it had a top speed of 32.5 knots, but "B" only required to make this speed. It also carried only nine 16-inch guns, in three triple turrets. "C" was similar but added (for a total of ) to meet the original requirement of . The weight required for this and a longer belt – , compared with for "B" – meant that the ship was standard.
Design history
In March 1938, the General Board followed the recommendations of the Battleship Design Advisory Board, which was composed of the naval architect William Francis Gibbs, William Hovgaard (then president of New York Shipbuilding), John Metten, Joseph W. Powell, and the long-retired Admiral and former Chief of the Bureau of Ordnance Joseph Strauss. The board requested an entirely new design study, again focusing on increasing the size of the South Dakota class. The first plans made for this indicated that was possible on a standard displacement of about . could be bought with and a standard displacement of around , which was well below the London Treaty's "escalator clause" maximum limit of .
These designs were able to convince the General Board that a reasonably well-designed and balanced 33-knot "fast" battleship was possible within the terms of the "escalator clause". However, further studies revealed major problems with the estimates. The speed of the ships meant that more freeboard would be needed both fore and amidships, the latter requiring an additional foot of armored freeboard. Along with this came the associated weight in supporting these new strains: the structure of the ship had to be reinforced and the power plant enlarged to avoid a drop in speed. In all, about had to be added, and the large margin the navy designers had previously thought they had – roughly – was suddenly vanishing. The draft of the ships was also allowed to increase, which enabled the beam to narrow and thus reduced the required power (since a lower beam-to-draft ratio reduces wave-making resistance). This also allowed the ships to be shortened, which reduced weight.
With the additional displacement, the General Board was incredulous that a tonnage increase of would allow only the addition of over the South Dakotas. Rather than retaining the 16-inch/45 caliber Mark 6 guns used in the South Dakotas, they ordered that the preliminary design would have to include the more powerful but significantly heavier 16-inch/50 caliber Mark 2 guns left over from the canceled s and battleships of the early 1920s.
The 16"/50 turret weighed some more than the 16"/45 turret already in use and also had a larger barbette diameter of compared to the latter's barbette diameter of , so the total weight gain was about . This put the ship at a total of – well over the limit. An apparent savior appeared in a Bureau of Ordnance preliminary design for a turret that could carry the 50-caliber guns and also fit in the smaller barbette of the 45-caliber gun turret. Other weight savings were achieved by thinning some armor elements and substituting construction steel with armor-grade Special Treatment Steel (STS) in certain areas. The net savings reduced the preliminary design displacement to standard, though the margin remained tight. This breakthrough was shown to the General Board as part of a series of designs on 2 June 1938.
However, the Bureau of Ordnance continued working on the turret with the larger barbette, while the Bureau of Construction and Repair used the smaller barbettes in the contract design of the new battleships. As the bureaus were independent of one another, they did not realize that the two plans could not go together until November 1938, when the contract design was in the final stages of refinement. By this time, the ships could not use the larger barbette, as it would require extensive alterations to the design and would result in substantial weight penalties. Reverting to the 45-caliber gun was also deemed unacceptable. The General Board was astounded; one member asked the head of the Bureau of Ordnance if it had occurred to him that Construction and Repair would have wanted to know what turret his subordinates were working on "as a matter of common sense". A complete scrapping of plans was avoided only when designers within the Bureau of Ordnance were able to design a new 50-caliber gun, the Mark 7, that was both lighter and smaller in outside diameter; this allowed it to be placed in a turret that would fit in the smaller barbette. The redesigned 3-gun turret, equipped as it was with the Mark 7 naval gun, provided an overall weight saving of nearly to the overall design of the Iowa class. The contract design displacement subsequently stood at standard and full load.
In May 1938, the United States Congress passed the Second Vinson Act, which "mandated a 20% increase in strength of the United States Navy". The act was sponsored by Carl Vinson, a Democratic Congressman from Georgia who was Chairman of the House Naval Affairs and Armed Services Committee. The Second Vinson Act updated the provisions of the Vinson-Trammell Act of 1934 and the Naval Act of 1936, which had "authorized the construction of the first American battleships in 17 years", based on the provisions of the London Naval Treaty of 1930; this act was quickly signed by President Franklin D. Roosevelt and provided the funding to build the Iowa class. Each ship cost approximately US$100 million.
As 1938 drew to a close, the contract design of the Iowas was nearly complete, but it would continuously evolve as the New York Navy Yard, the lead shipyard, conducted the final detail design. These revisions included changing the design of the foremast, replacing the original /75-caliber guns that were to be used for anti-aircraft (AA) work with /70 caliber Oerlikon cannons and /56 caliber Bofors guns, and moving the combat information center into the armored hull. Additionally, in November 1939, the New York Navy Yard greatly modified the internal subdivision of the machinery rooms, as tests had shown the underwater protection in these rooms to be inadequate. The longitudinal subdivision of these rooms was doubled, and the result of this was clearly beneficial: "The prospective effect of flooding was roughly halved and the number of uptakes and hence of openings in the third deck greatly reduced." Although the changes meant extra weight and increasing the beam by to , this was no longer a major issue; Britain and France had renounced the Second London Naval Treaty soon after the beginning of the Second World War. The design displacement was standard, approximately 2% overweight, when Iowa and New Jersey were laid down in June and September 1940. By the time the Iowas were completed and commissioned in 1943–44, the considerable increase in anti-aircraft armament – along with their associated splinter protection and crew accommodations – and additional electronics had increased standard displacement to some , while full load displacement became .
Specifications
General characteristics
The Iowa-class battleships are long at the waterline and long overall with a beam of . During World War II, the draft was at full load displacement of and at design combat displacement of . Like the two previous classes of American fast battleships, the Iowas have a double bottom hull that becomes a triple bottom under the armored citadel and armored skegs around the inboard shafts. The dimensions of the Iowas were strongly influenced by speed. When the Second Vinson Act was passed by the United States Congress in 1938, the U.S. Navy moved quickly to develop a 45,000-ton-standard battleship that would pass through the wide Panama Canal. Drawing on a 1935 empirical formula for predicting a ship's maximum speed based on scale-model studies in flumes of various hull forms and propellers and a newly developed empirical theorem that related waterline length to maximum beam, the Navy drafted plans for a battleship class with a maximum beam of which, when multiplied by 7.96, produced a waterline length of . The Navy also called for the class to have a lengthened forecastle and amidship, which would increase speed, and a bulbous bow.
The Iowas exhibit good stability, making them steady gun platforms. At design combat displacement, the ships' (GM) metacentric height was . They also have excellent maneuverability in the open water for their size, while seakeeping is described as good, but not outstanding. In particular, the long fine bow and sudden widening of the hull just in front of the foremost turret contributed to the ships being rather wet for their size. This hull form also resulted in very intense spray formations, which led to some difficulty refueling escorting destroyers.
Armament
Main battery
The primary guns used on these battleships are the nine /50-caliber Mark 7 naval guns, a compromise design developed to fit inside the barbettes. These guns fire high explosive- and armor-piercing shells and can fire a 16-inch shell approximately . The guns are housed in three 3-gun turrets: two forward of the battleship's superstructure and one aft, in a configuration known as "2-A-1". The guns are long (50 times their 16-inch bore, or 50 calibers from breechface to muzzle). About protrudes from the gun house. Each gun weighs about without the breech, or with the breech. They fired armor-piercing projectiles at a muzzle velocity of , or high-capacity projectiles at , up to . At maximum range, the projectile spends almost minutes in flight. The maximum firing rate for each gun is two rounds per minute.
Each gun rests within an armored turret, but only the top of the turret protrudes above the main deck. The turret extends either four decks (Turrets 1 and 3) or five decks (Turret 2) down. The lower spaces contain rooms for handling the projectiles and storing the powder bags used to fire them. Each turret required a crew of between 85 and 110 men to operate. The original cost for each turret was US$1.4 million, but this figure does not take into account the cost of the guns themselves. The turrets are "three-gun", not "triple", because each barrel is individually sleeved and can be elevated and fired independently. The ship could fire any combination of its guns, including a broadside of all nine. The fire control was performed by the Mark 38 Gun Fire Control System (GFCS); the firing solutions were computed with the Mark 8 rangekeeper, an analog computer that automatically receives information from the director and Mark 8/13 fire control radar, stable vertical, ship pitometer log and gyrocompass, and anemometer. The GFCS uses remote power control (RPC) for automatic gun laying.
The large-caliber guns were designed to fire two different conventional 16-inch shells: the Mk 8 "Super-heavy" APC (Armor Piercing, Capped) shell for anti-ship and anti-structure work, and the Mk 13 high-explosive round designed for use against unarmored targets and shore bombardment. When firing the same conventional shell, the 16-inch/45 caliber Mark 6 used by the fast battleships of the North Carolina and South Dakota classes had a slight advantage over the 16-inch/50 caliber Mark 7 gun when hitting deck armor – a shell from a 45 cal gun would be slower, meaning that it would have a steeper trajectory as it descended. At , a shell from a 45 cal would strike a ship at an angle of 45.2 degrees, as opposed to 36 degrees with the 50 cal. The Mark 7 had a greater maximum range over the Mark 6: vs .
In the 1950s, the W23, an adaptation of the W19 nuclear artillery shell, was developed specifically for the 16-inch guns. The shell weighed , had an estimated yield of , and its introduction made the Iowa-class battleships' 16-inch guns the world's largest nuclear artillery and made these four battleships the only US Navy ships ever to have nuclear shells for naval guns. Although developed for exclusive use by the battleship's guns it is not known if any of the Iowas actually carried these shells while in active service due to the United States Navy's policy of refusing to confirm or deny the presence of nuclear weaponry aboard its ships. In 1991, the United States unilaterally withdrew all of its nuclear artillery shells from service, and the dismantling of the US nuclear artillery inventory is said to have been completed in 2004.
Secondary battery
The Iowas carried twenty /38 caliber Mark 12 guns in ten Mark 28 Mod 2 enclosed base ring mounts. Originally designed to be mounted upon destroyers built in the 1930s, these guns were so successful that they were added to many American ships during the Second World War, including every major ship type and many smaller warships constructed between 1934 and 1945. They were considered to be "highly reliable, robust and accurate" by the Navy's Bureau of Ordnance.
Each 5-inch/38 gun weighed almost without the breech; the entire mount weighed . It was long overall, had a bore length of , and a rifling length of . The gun could fire shells at about ; about 4,600 could be fired before the barrel needed to be replaced. Minimum and maximum elevations were −15 and 85 degrees, respectively. The guns' elevation could be raised or lowered at about 15 degrees per second. The mounts closest to the bow and stern could aim from −150 to 150 degrees; the others were restricted to −80 to 80 degrees. They could be turned at about 25 degrees per second. The mounts were directed by four Mark 37 fire control systems primarily through remote power control (RPC).
The 5-inch/38 gun functioned as a dual-purpose gun (DP); that is, it was able to fire at both surface and air targets with a reasonable degree of success. However, this did not mean that it possessed inferior anti-air abilities. As proven during 1941 gunnery tests conducted aboard the gun could consistently shoot down aircraft flying at , twice the effective range of the earlier single-purpose 5-inch/25 caliber AA gun. As Japanese airplanes became faster, the gun lost some of its effectiveness in the anti-aircraft role; however, toward the end of the war, its usefulness as an anti-aircraft weapon increased again because of an upgrade to the Mark 37 Fire Control System, Mark 1A computer, and proximity-fused shells.
The 5-inch/38 gun would remain on the battleships for the ships' entire service life; however, the total number of guns and gun mounts was reduced from twenty guns in ten mounts to twelve guns in six mounts during the 1980s' modernization of the four Iowas. The removal of four of the gun mounts was required for the battleships to be outfitted with the armored box launchers needed to carry and fire Tomahawk missiles. At the time of the 1991 Persian Gulf War, these guns had been largely relegated to littoral defense for the battleships. Since each battleship carried a small detachment of Marines aboard, the Marines would man one of the 5-inch gun mounts.
Anti-air battery
At the time of their commissioning, all four of the Iowa-class battleships were equipped with 20 quad 40 mm mounts and 49 single 20 mm mounts. These guns were respectively augmented with the Mk 14 range sight and Mk 51 fire control system to improve accuracy.
The Oerlikon gun, one of the most heavily produced anti-aircraft guns of the Second World War, entered service in 1941 and replaced the M2 Browning MG on a one-for-one basis. Between December 1941 and September 1944, 32% of all Japanese aircraft downed were credited to this weapon, with the high point being 48.3% for the second half of 1942; however, the 20 mm guns were found to be ineffective against the Japanese Kamikaze attacks used during the latter half of World War II and were subsequently phased out in favor of the heavier Bofors AA gun.
When the Iowa-class battleships were commissioned in 1943 and 1944, they carried twenty quad 40 mm AA gun mounts, which they used for defense against enemy aircraft. These heavy AA guns were also employed in the protection of Allied aircraft carriers operating in the Pacific Theater of World War II, and accounted for roughly half of all Japanese aircraft shot down between 1 October 1944 and 1 February 1945. Although successful in this role against WWII aircraft, the 40 mm guns were stripped from the battleships in the jet age – initially from New Jersey when reactivated in 1968 and later from Iowa, Missouri, and Wisconsin when they were reactivated for service in the 1980s.
Propulsion
The powerplant of the Iowas consists of eight Babcock & Wilcox boilers and four sets of double reduction cross-compound geared turbines, with each turbine set driving a single shaft. Specifically, the geared turbines on Iowa and Missouri were provided by General Electric, while the equivalent machinery on New Jersey and Wisconsin was provided by Westinghouse. The plant produced and propelled the ship up to a maximum speed of at full load displacement and at normal displacement. The ships carried of fuel oil which gave a range of at . Two semi-balanced rudders gave the ships a tactical turning diameter of at and at .
The machinery spaces were longitudinally divided into eight compartments with alternating fire and engine rooms to ensure adequate isolation of machinery components. Four fire rooms each contained two M-Type boilers operating at with a maximum superheater outlet temperature of . The double-expansion engines consist of a high-pressure (HP) turbine and a low-pressure (LP) turbine. The steam is first passed through the HP turbine which turns at up to 2,100 rpm. The steam, largely depleted at this point, is then passed through a large conduit to the LP turbine. By the time it reaches the LP turbine, it has no more than of pressure left. The LP turbine increases efficiency and power by extracting the last little bit of energy from the steam. After leaving the LP turbine, the exhaust steam passes into a condenser and is then returned as feed water to the boilers. Water lost in the process is replaced by three evaporators, which can make a total of 60,000 US gallons per day (3 liters per second) of fresh water. After the boilers have had their fill, the remaining fresh water is fed to the ship's potable water systems for drinking, showers, hand washing, cooking, etc. All of the urinals and all but one of the toilets on the Iowa class flush with salt water in order to conserve fresh water. The turbines, especially the HP turbine, can turn at 2,000 rpm; their shafts drive through reduction gearing that turns the propeller shafts at speeds up to 225 rpm, depending upon the desired speed of the ship. The Iowas were outfitted with four screws: the outboard pair consisting of four-bladed propellers in diameter and the inboard pair consisting of five-bladed propellers in diameter. The propeller designs were adopted after earlier testing had determined that propeller cavitation caused a drop in efficiency at speeds over . The two inner shafts were housed in skegs to smooth the flow of water to the propellers and improve the structural strength of the stern.
Each of the four engine rooms has a pair of 1,250 kW Ship's Service Turbine Generators (SSTGs), providing the ship with a total non-emergency electrical power of 10,000 kW at 450 volts alternating current. Additionally, the vessels have a pair of 250 kW emergency diesel generators. To allow battle-damaged electrical circuits to be repaired or bypassed, the lower decks of the ship have a Casualty Power System whose large 3-wire cables and wall outlets called "biscuits" can be used to reroute power.
Electronics (1943–69)
The earliest search radars installed were the SK air-search radar and SG surface-search radar during World War II. They were located on the mainmast and forward fire-control tower of the battleships, respectively. As the war drew to a close, the United States introduced the SK-2 air-search radar and SG surface-search radar; the Iowa class was updated to make use of these systems between 1945 and 1952. At the same time, the ships' radar systems were augmented with the installation of the SP height finder on the main mast. In 1952, AN/SPS-10 surface-search radar and AN/SPS-6 air-search radar replaced the SK and SG radar systems, respectively. Two years later the SP height finder was replaced by the AN/SPS-8 height finder, which was installed on the main mast of the battleships.
In addition to these search and navigational radars, the Iowa class were also outfitted with a variety of fire control radars for their gun systems. Beginning with their commissioning, the battleships made use of a pair of Mk 38 gun fire control systems with Mark 8 fire control radar to direct the 16-inch guns and a quartet of Mk 37 gun fire control systems with Mark 12 fire control radar and Mark 22 height finding radar to direct the 5-inch gun batteries. These systems were upgraded over time with the Mark 13 replacing the Mark 8 and the Mark 25 replacing the Mark 12/22, but they remained the cornerstones of the combat radar systems on the Iowa class during their careers. The range estimation of these gunfire control systems provided a significant accuracy advantage over earlier ships with optical rangefinders; this was demonstrated off Truk Atoll on 16 February 1944, when the New Jersey engaged the at a range of and straddled her, setting the record for the longest-ranged straddle in history.
In World War II, the electronic countermeasures (ECM) included the SPT-1 and SPT-4 equipment; passive electronic support measures (ESM) were a pair of DBM radar direction finders and three intercept receiving antennas, while the active components were the TDY-1 jammers located on the sides of the fire control tower. The ships were also equipped with the identification, friend or foe (IFF) Mark III system, which was replaced by the IFF Mark X when the ships were overhauled in 1955. When the New Jersey was reactivated in 1968 for the Vietnam War, she was outfitted with the ULQ-6 ECM system.
Armor
Like all battleships, the Iowas carried heavy armor protection against shellfire and bombs with significant underwater protection against torpedoes. The Iowas' "all-or-nothing" armor scheme was largely modeled on that of the preceding South Dakota class, and designed to give a zone of immunity against fire from 16-inch/45-caliber guns between away. The protection system consists of Class A face-hardened Krupp cemented (K.C.) armor and Class B homogeneous Krupp-type armor; furthermore, special treatment steel (STS), a high-tensile structural steel with armor properties comparable to Class B, was extensively used in the hull plating to increase protection.
The citadel consisting of the magazines and engine rooms was protected by an STS outer hull plating thick and a Class A armor belt thick mounted on STS backing plate; the armor belt is sloped at 19 degrees, equivalent to of vertical class B armor at 19,000 yards. The armor belt extends to the triple bottom, where the Class B lower portion tapers to . The ends of the armored citadel are closed by vertical Class A transverse bulkheads for Iowa and New Jersey. The transverse bulkhead armor on Missouri and Wisconsin was increased to ; this extra armor provided protection from raking fire directly ahead, which was considered more likely given the high speed of the Iowas. The deck armor consists of a STS weather deck, a combined Class B and STS main armor deck, and a STS splinter deck. Over the magazines, the splinter deck is replaced by a STS third deck that separates the magazine from the main armored deck. The powder magazine rooms are separated from the turret platforms by a pair of 1.5-inch STS annular bulkheads under the barbettes for flashback protection. The installation of armor on the Iowas also differed from those of earlier battleships in that the armor was installed while the ships were still "on the way" rather than after the ships had been launched.
The Iowas had heavily protected main battery turrets, with Class B and STS face, Class A sides, Class A rear, and Class B roof. The turret barbettes' armor is Class A with abeam and facing the centerline, extending down to the main armor deck. The conning tower armor is Class B with on all sides and on the roof. The secondary battery turrets and handling spaces were protected by of STS. The propulsion shafts and steering gear compartment behind the citadel had considerable protection, with Class A side strake and roof.
The armor's immunity zone shrank considerably against guns equivalent to their own 16-inch/50-caliber guns armed with the Mk 8 armor-piercing shell due to the weapon's increased muzzle velocity and improved shell penetration; increasing the armor would have increased weight and reduced speed, a compromise that the General Board was not willing to make.
The Iowas' torpedo defense was based on the South Dakotas' design, with modifications to address shortcomings discovered during caisson tests. The system is an internal "bulge" that consists of four longitudinal torpedo bulkheads behind the outer hull plating with a system depth of to absorb the energy of a torpedo warhead. The extension of the armor belt to the triple bottom, where it tapers to a thickness of , serves as one of the torpedo bulkheads and was hoped to add to protection; the belt's lower edge was welded to the triple bottom structure and the joint was reinforced with buttstraps due to the slight knuckle causing a structural discontinuity. The torpedo bulkheads were designed to elastically deform to absorb energy and the two outer compartments were liquid loaded in order to disrupt the gas bubble and slow fragments. The outer hull was intended to detonate a torpedo, with the outer two liquid compartments absorbing the shock and slowing any splinters or debris while the lower armored belt and the empty compartment behind it absorb any remaining energy. However, the Navy discovered in caisson tests in 1939 that the initial design for this torpedo defense system was actually less effective than the previous design used on the North Carolinas due to the rigidity of the lower armor belt causing the explosion to significantly displace the final holding bulkhead inwards despite remaining watertight. To mitigate the effects, the third deck and triple bottom structure behind the lower armor belt were reinforced and the placement of brackets was changed. Iowas' system was also improved over the South Dakotas' through closer spacing of the transverse bulkheads, greater thickness of the lower belt at the triple bottom joint, and increased total volume of the "bulge". The system was further modified for the last two ships of the class, Illinois and Kentucky, by eliminating knuckles along certain bulkheads; this was estimated to improve the strength of the system by as much as 20%.
Based on costly lessons in the Pacific theater, concerns were raised about the ability of the armor on these battleships to withstand aerial bombing, particularly high-altitude bombing using armor-piercing bombs. Developments such as the Norden bombsight further fueled these concerns. While the design of the Iowas was too far along to adequately address this issue, experience in the Pacific theater eventually demonstrated that high-altitude unguided bombing was ineffective against maneuvering warships.
Aircraft (1943–69)
When they were commissioned during World War II, the Iowa-class battleships came equipped with two aircraft catapults designed to launch floatplanes. Initially, the Iowas carried the Vought OS2U Kingfisher and Curtiss SC Seahawk, both of which were employed to spot for the battleship's main gun batteries – and, in a secondary capacity, perform search-and-rescue missions.
By the time of the Korean War, helicopters had replaced floatplanes and the Sikorsky HO3S-1 helicopter was employed. New Jersey made use of the Gyrodyne QH-50 DASH drone for her Vietnam War deployment in 1968–69.
Conversion proposals
The Iowa class were the only battleships with the speed required for post-war operations based around fast aircraft carrier task forces. There were several proposals in the early Cold War to convert the class to take into account changes in technology and doctrine. These included plans to equip the class with nuclear missiles, add aircraft capability, and – in the case of Illinois and Kentucky – a proposal to rebuild both as aircraft carriers instead of battleships.
Initially, the Iowa class was to consist of only four battleships with hull numbers BB-61 to BB-64: Iowa, New Jersey, Missouri, and Wisconsin. However, changing priorities during World War II resulted in the battleship hull numbers BB-65 and BB-66 being reordered as Illinois and Kentucky, respectively; Montana and Ohio were reassigned to hull numbers BB-67 and BB-68. At the time these two battleships were to be built a proposal was put forth to have them constructed as aircraft carriers rather than fast battleships. The plan called for the ships to be rebuilt to include a flight deck and an armament suite similar to that placed aboard the s that were at the time under construction in the United States. Ultimately, nothing came of the design proposal to rebuild these two ships as aircraft carriers and they were cleared for construction as fast battleships to conform to the Iowa-class design, though they differed from the earlier four that were built. Eventually, the light cruisers were selected for the aircraft-carrier conversion. Nine of these light cruisers would be rebuilt as light aircraft carriers.
After the surrender of the Empire of Japan, construction on Illinois and Kentucky stopped. Illinois was eventually scrapped, but Kentuckys construction had advanced enough that several plans were proposed to complete Kentucky as a guided missile battleship (BBG) by removing the aft turret and installing a missile system. A similar conversion had already been performed on the battleship (BB-41/AG-128) to test the RIM-2 Terrier missile after World War II. One such proposal came from Rear Admiral W.K. Mendenhall, Chairman of the Ship Characteristics Board (SCB); Mendenhall proposed a plan that called for $15–30 million to be spent to allow Kentucky to be completed as a guided-missile battleship (BBG) carrying eight SSM-N-8 Regulus II guided missiles with a range of . He also suggested Terrier or RIM-8 Talos launchers to supplement the AA guns and proposed nuclear (instead of conventional) shells for the 16-inch guns. This never materialized, and Kentucky was ultimately sold for scrap in 1958, although her bow was used to repair her sister Wisconsin after a collision on 6 May 1956, earning her the nickname WisKy.
In 1954, the Long Range Objectives Group of the United States Navy suggested converting the Iowa-class ships to BBGs. In 1958, the Bureau of Ships offered a proposal based on this idea. This replaced the 5- and 16-inch gun batteries with "two Talos twin missile systems, two RIM-24 Tartar twin missile systems, an RUR-5 ASROC antisubmarine missile launcher, and a Regulus II installation with four missiles", as well as flagship facilities, sonar, helicopters, and fire-control systems for the Talos and Tartar missiles. In addition to these upgrades, of additional fuel oil was also suggested to serve in part as ballast for the battleships and for use in refueling destroyers and cruisers. Due to the estimated cost of the overhaul ($178–193 million) this proposal was rejected as too expensive; instead, the SCB suggested a design with one Talos, one Tartar, one ASROC, and two Regulus launchers and changes to the superstructure, at a cost of up to $85 million. This design was later revised to accommodate the Polaris Fleet Ballistic Missile, which in turn resulted in a study of two schemes by the SCB. In the end, none of these proposed conversions for the battleships were ever authorized. Interest in converting the Iowas into guided-missile battleships began to deteriorate in 1960 because the hulls were considered too old and the conversion costs too high. Nonetheless, additional conversion proposals – including one to install the AN/SPY-1 Aegis Combat System radar on the battleships – were suggested in 1962, 1974, and 1977, but as before, these proposals failed to gain the needed authorization. This was due, in part, to the possibility that sensitive electronics within of any 16-inch gun muzzle may be damaged from overpressure.
1980s reactivation/modernization
In 1980, Ronald Reagan was elected president on a promise to build up the US military as a response to the increasing military power of the Soviet Union. The Soviet Navy was commissioning the Kirov class of missile cruisers, the largest type of surface combatant since World War II. As part of Reagan's 600-ship Navy policy and as a counter to the Kirov class, the US Navy began reactivating the four Iowa-class units and modernizing them for service.
The Navy considered several proposals that would have removed the aft 16-inch turret. Martin Marietta proposed to replace the turret with servicing facilities for 12 AV-8B Harrier STOVL jump jets. A more detailed proposal, the "Interdiction Assault Ship", proposed a V-shaped ramped flight deck (the base of the V would have been on the ship's stern, while each leg of the V would extend forward, so that planes taking off would fly past the ship's exhaust stacks and conning tower), while a new hangar would be added with two elevators, which would support up to twelve McDonnell Douglas AV-8B Harrier II jump-jets. These aviation facilities could also support helicopters, SEAL teams and up to 500 Marines for an air assault. In the empty space between the V flight deck would be up to 320 missile silos accommodating a mixture of Tomahawk land attack missiles, ASROC anti-submarine rockets and Standard surface-to-air missiles. The existing five-inch gun turrets would be replaced with 155-millimeter howitzers for naval gunfire support.
Charles Myers, a former Navy test pilot turned Pentagon consultant, proposed replacing the turret with vertical launch systems for missiles and a flight deck for Marine helicopters. In July 1981, the US Naval Institute's Proceedings published a proposal by naval architect Gene Anderson for a canted flight deck with steam catapult and arrestor wires for F/A-18 Hornet fighters. Plans for these conversions were dropped in 1984.
Each battleship was overhauled to burn navy distillate fuel and modernized to carry electronic warfare suites, close-in weapon systems (CIWS) for self-defense, and missiles. The obsolete electronics and anti-aircraft armament were removed to make room for more modern systems. The Navy spent about $1.7 billion, from 1981 through 1988, to modernize and reactivate the four Iowa-class battleships, roughly the same as building four Oliver Hazard Perry-class frigates.
After modernization, the full load displacement was relatively unchanged at .
The modernized battleships operated as centerpieces of their own battle group (termed as a Battleship Battle Group or Surface Action Group), consisting of one , one or , one , three s and one support ship, such as a fleet oiler.
Armament
During their modernization in the 1980s, each Iowa was equipped with four of the US Navy's Phalanx CIWS mounts, two of which sat just behind the bridge and two which were next to the ship's aft funnel. Iowa, New Jersey, and Missouri were equipped with the Block 0 version of the Phalanx, while Wisconsin received the first operational Block 1 version in 1988. The Phalanx system is intended to serve as a last line of defense against enemy missiles and aircraft, and when activated can engage a target with a 20 mm M61 Vulcan 6-barreled Gatling cannon at a distance of approximately .
As part of their modernization in the 1980s, each of the Iowas received a complement of eight quad-cell Armored Box Launchers and four "shock hardened" Mk 141 quad-cell launchers. The former was used by the battleships to carry and fire the BGM-109 Tomahawk Land Attack Missiles (TLAMs) for use against enemy targets on land, while the latter system enabled the ships to carry a complement of RGM-84 Harpoon anti-ship missiles for use against enemy ships. With an estimated range of for the Tomahawks and for the Harpoons, these two missile systems displaced the 16-inch guns and their maximum range of to become the longest-ranged weapons on the battleships during the 1980s; the ships' complement of 32 Tomahawk missiles was the largest until the Mk 41 VLS-equipped cruisers entered service. It has been alleged by members of the environmental group Greenpeace that the battleships carried the TLAM-A (also cited, incorrectly, as the TLAM-N) – a Tomahawk missile with a variable yield W80 nuclear warhead – during their 1980s service with the United States Navy, but owing to the United States Navy's policy of refusing to confirm or deny the presence of nuclear weaponry aboard its ships, these claims can not be conclusively proved. Between 2010 and 2013, the US withdrew the BGM-109A, leaving only conventional munitions packages for its Tomahawk missile inventory, though the Iowas had been withdrawn from service at that point.
Owing to the original 1938 design of the battleships, the Tomahawk missiles could not be fitted to the Iowa class unless the battleships were rebuilt in such a way as to accommodate the missile mounts that would be needed to store and launch the Tomahawks. This realization prompted the removal of the anti-aircraft guns previously installed on the Iowas and the removal of four of each of the battleships' ten 5-inch/38 DP mounts. The mid and aft end of the battleships were then rebuilt to accommodate the missile launchers. At one point, the NATO Sea Sparrow was to be installed on the reactivated battleships; however, it was determined that the system could not withstand the overpressure effects of firing the main battery. To supplement the anti-aircraft capabilities of the Iowas, five FIM-92 Stinger surface-to-air missile firing positions were installed. These secured the shoulder-launched weapons and their rounds for ready use by the crew.
Electronics
During their modernization under the 600-ship Navy program, the Iowa-class battleships' radar systems were again upgraded. The foremast was of a new tripod design that was considerably reinforced to allow the AN/SPS-6 air-search radar system to be replaced with the AN/SPS-49 radar set (which also augmented the existing navigation capabilities on the battleships), and the AN/SPS-8 surface-search radar set was replaced by the AN/SPS-67 search radar. The new mast also incorporates a Tactical Air Navigation System (TACAN) antenna. The aft mast was changed to be placed in front of the aft funnel and mounts a circular SATCOM antenna while another one was mounted on the fire control mast.
By the Korean War, jet engines had replaced propellers on aircraft, which severely limited the ability of the 20 mm and 40 mm AA batteries and their gun systems to track and shoot down enemy planes. Consequently, the AA guns and their associated fire-control systems were removed when reactivated. New Jersey received this treatment in 1967, and the others followed in their 1980s modernizations. In the 1980s, each ship also received a quartet of Phalanx CIWS mounts which made use of a radar system to locate incoming enemy projectiles and destroy them with a 20 mm Gatling gun before they could strike the ship.
With the added missile capacity of the battleships in the 1980s came additional fire-support systems to launch and guide the ordnance. To fire the Harpoon anti-ship missiles, the battleships were equipped with the SWG-1 fire-control system, and to fire the Tomahawk missiles the battleships used either the SWG-2 or SWG-3 fire-control system. In addition to these offensive-weapon systems, the battleships were outfitted with the AN/SLQ-25 Nixie to be used as a lure against enemy torpedoes; an SLQ-32 electronic warfare system that can detect, jam, and deceive an opponent's radar; and a Mark 36 SRBOC system to fire chaff rockets intended to confuse enemy missiles.
Aside from the electronics added for weaponry control, all four battleships were outfitted with a communications suite used by both cruisers and guided missile cruisers in service at the time. This communication suite included the OE-82 antenna for satellite communications but did not include the Naval Tactical Data System.
Aircraft (1982–1992)
During the 1980s these battleships made use of the RQ-2 Pioneer, an unmanned aerial vehicle employed in spotting for the guns. Launched from the fantail using a rocket-assist booster that was discarded shortly after takeoff, the Pioneer carried a video camera in a pod under the belly of the aircraft which transmitted live video to the ship so operators could observe enemy actions or fall of shot during naval gunnery. To land the UAV a large net was deployed at the back of the ship; the aircraft was flown into it. Missouri and Wisconsin both used the Pioneer UAVs successfully during Operation Desert Storm, and in one particularly memorable incident, a Pioneer UAV operated by Wisconsin received the surrender of Iraqi troops during combat operations. This particular Pioneer was later donated to the Smithsonian Institution and is now on public display. During Operation Desert Storm these Pioneers were operated by detachments of VC-6. In addition to the Pioneer UAVs, the recommissioned Iowas could support operations by various types of helicopters, including the UH-1 Iroquois, SH-2 Seasprites, CH-46 Sea Knight, CH-53 Sea Stallion, and LAMPS III SH-60B Seahawk.
Gunfire support role
Following the 1991 Gulf War and the subsequent dissolution of the Soviet Union, the United States Navy began to decommission and mothball many of the ships it had brought out of its reserve fleet in the drive to attain a 600-ship Navy. At the height of Navy Secretary John F. Lehman's 600-ship Navy plan, nearly 600 ships of all types were active within the Navy. This included fifteen aircraft carriers, four battleships, and over 100 submarines, along with various other types of ships the overall plan specified. When the Soviet Union collapsed in 1991 the Navy sought to return to its traditional, 313-ship composition. While reducing the fleet created under the 600-ship Navy program, the decision was made to deactivate the four recommissioned Iowa-class battleships and return them to the reserve fleet.
In 1995, the decommissioned battleships were removed from the Naval Vessel Register after it was determined by ranking US Navy officials that there was no place for a battleship in the modern navy. In response to the striking of the battleships from the Naval Vessel Register a movement began to reinstate the battleships, on the grounds that these vessels had superior firepower over the 5-inch guns found on the Spruance, Kidd and Arleigh Burke-class destroyers and Ticonderoga-class cruisers. Citing concern over the lack of available gunfire to support amphibious operations, Congress required the Navy to reinstate two battleships to the Naval Vessel Register and maintain them with the mothball fleet, until the Navy could certify it had gunfire support within the current fleet that would meet or exceed the battleship's capability.
The debate over battleships in the modern navy continued until 2006, when the two reinstated battleships were stricken after naval officials submitted a two-part plan that called for the near-term goal of increasing the range of the guns in use on the Arleigh Burke-class destroyers with new Extended Range Guided Munition (ERGM) ammunition intended to allow a 5-inch projectile fired from these guns to travel an estimated inland. The long-term goal called for the replacement of the two battleships with 32 vessels of the of guided-missile destroyers. Cost overruns caused the class to be reduced to three ships. These ships are outfitted with an Advanced Gun System (AGS) that was to fire specially developed 6-inch Long Range Land Attack Projectiles for shore bombardment. LRLAP procurement was canceled in 2017 and the AGS is unusable. The long-term goal for the Zumwalt class is to have the ships mount railguns or free-electron lasers.
Cultural significance
The Iowa class became culturally symbolic in the United States in many different ways, to the point where certain elements of the American public – such as the United States Naval Fire Support Association – were unwilling to part with the battleships, despite their apparent obsolescence in the face of modern naval combat doctrine that places great emphasis on air supremacy and missile firepower. Although all were officially stricken from the Naval Vessel Register they were spared scrapping and were donated for use as museum ships.
Their service records added to their fame, ranging from their work as carrier escorts in World War II to their shore bombardment duties in North Korea, North Vietnam, and the Middle East, as well as their service in the Cold War against the expanded Soviet Navy. Their reputation combined with the stories told concerning the firepower of these battleships' 16-inch guns were such that when they were brought out of retirement in the 1980s in response to increased Soviet Naval activity – and in particular, in response to the commissioning of the Kirov-class battlecruisers – the United States Navy was inundated with requests from former sailors pleading for a recall to active duty so they could serve aboard one of the battleships.
In part because of the service length and record of the class, members have made numerous appearances in television shows, video games, movies, and other media, including appearances of the Kentucky and Illinois in the anime series Neon Genesis Evangelion, the History Channel documentary series Battle 360: USS Enterprise, the Discovery Channel documentary The Top 10 Fighting Ships (where the Iowa class was rated Number 1), the book turned movie A Glimpse of Hell, the 1989 music video for the song by Cher "If I Could Turn Back Time", the 1992 film Under Siege, the 2012 film Battleship, among other appearances. Japanese rock band Vamps performed the finale of their 2009 US tour on board Missouri on 19 September 2009.
Ships in class
When brought into service during the final years of World War II, the Iowa-class battleships were assigned to operate in the Pacific Theatre of World War II. By this point in the war, aircraft carriers had displaced battleships as the primary striking arm of both the United States Navy and the Imperial Japanese Navy. As a result of this shift in tactics, US fast battleships of all classes were relegated to the secondary role of carrier escorts and assigned to the Fast Carrier Task Force to provide anti-aircraft screening for Allied aircraft carriers and perform shore bombardment. Three were recalled to service in the 1950s with the outbreak of the Korean War, and they provided naval artillery support for U.N. forces for the entire duration of the war before being returned to mothballs in 1955 after hostilities ceased. In 1968, to help alleviate US air losses over North Vietnam, New Jersey was summoned to Vietnam, but she was decommissioned a year after arriving. All four returned in the 1980s during the drive for a 600-ship Navy to counter the new Soviet Kirov-class battlecruisers, only to be retired after the collapse of the Soviet Union on the grounds that they were too expensive to maintain.
Iowa
Iowa was ordered 1 July 1939, laid down 27 June 1940, launched 27 August 1942, and commissioned 22 February 1943. She conducted a shakedown cruise in Chesapeake Bay before sailing to Naval Station Argentia, Newfoundland, to be ready in case the entered the Atlantic. Transferred to the Pacific Fleet in 1944, Iowa made her combat debut in February and participated in the campaign for the Marshall Islands. The ship later escorted US aircraft carriers conducting air raids in the Marianas campaign, and then was present at the Battle of Leyte Gulf. During the Korean War, Iowa bombarded enemy targets at Songjin, Hŭngnam and Kojo, North Korea. Iowa returned to the US for operational and training exercises before being decommissioned on 24 February 1958. Reactivated in the early 1980s, Iowa operated in the Atlantic Fleet, cruising in North American and European waters for most of the decade and participating in joint military exercises with European ships. On 19 April 1989, 47 sailors were killed following an explosion in her No. 2 turret. In 1990, Iowa was decommissioned for the last time and placed in the mothball fleet. She was stricken from the Naval Vessel Register on 17 March 2006. Iowa was anchored as part of the National Defense Reserve Fleet in Suisun Bay, California until October 2011, when she was towed from her mooring to Richmond, California for renovation as a museum ship. She was towed from Richmond in the San Francisco Bay on 26 May 2012, to San Pedro at the Los Angeles Waterfront to serve as a museum ship run by Pacific Battleship Center and opened to the public on 7 July 2012.
New Jersey
New Jersey was ordered 4 July 1939, laid down 16 September 1940, launched 7 December 1942, and commissioned 23 May 1943. New Jersey completed fitting out and trained her initial crew in the Western Atlantic and Caribbean before transferring to the Pacific Theatre in advance of the planned assault on the Marshall Islands, where she screened the US fleet of aircraft carriers from enemy air raids. At the Battle of Leyte Gulf, the ship protected carriers with her anti-aircraft guns. New Jersey then bombarded Iwo Jima and Okinawa. During the Korean War, the ship pounded targets at Wonsan, Yangyang, and Kansong. Following the Armistice, New Jersey conducted training and operation cruises until she was decommissioned on 21 August 1957. Recalled to duty in 1968, New Jersey reported to the gunline off the Vietnamese coast and shelled North Vietnamese targets before departing the line in December 1968. She was decommissioned the following year. Reactivated in 1982 under the 600-ship Navy program, New Jersey was sent to Lebanon to protect US interests and US Marines, firing her main guns at Druze and Syrian positions in the Beqaa Valley east of Beirut. Decommissioned for the last time 8 February 1991, New Jersey was briefly retained on the Naval Vessel Register before being donated to the Home Port Alliance of Camden, New Jersey for use as a museum ship in October 2001.
Missouri
Missouri was the last of the four Iowas to be completed. She was ordered 12 June 1940, laid down 6 January 1941, launched 29 January 1944, and commissioned 11 June 1944. Missouri conducted her trials off New York with shakedown and battle practice in the Chesapeake Bay before transferring to the Pacific Fleet, where she screened US aircraft carriers involved in offensive operations against the Japanese before reporting to Okinawa to shell the island in advance of the planned landings. Following the bombardment of Okinawa, Missouri turned her attention to the Japanese homeland islands of Honshu and Hokkaido, performing shore bombardment and screening US carriers involved in combat operations. She became a symbol of the US Navy's victory in the Pacific when representatives of the Empire of Japan boarded the battleship to sign the documents of unconditional surrender to the Allied powers in September 1945. After World War II, Missouri conducted largely uneventful training and operational cruises until suffering a grounding accident. In 1950, she was dispatched to Korea in response to the outbreak of the Korean War. Missouri served two tours of duty in Korea providing shore bombardment. She was decommissioned in 1956. She spent many years at Puget Sound Naval Shipyard in Bremerton, Washington. Reactivated in 1984, as part of the 600-ship Navy plan, Missouri was sent on operational cruises until being assigned to Operation Earnest Will in 1988. In 1991, Missouri participated in Operation Desert Storm, firing 28 Tomahawk Missiles and 759 16-inch shells at Iraqi targets along the coast. Decommissioned for the last time in 1992, Missouri was donated to the USS Missouri Memorial Association of Pearl Harbor, Hawaii, for use as a museum ship in 1999.
Wisconsin
Wisconsin was ordered 12 June 1940, laid down 25 January 1942, launched 7 December 1943, and commissioned 16 April 1944. After trials and initial training in the Chesapeake Bay, she transferred to the Pacific Fleet in 1944 and was assigned to protect the US fleet of aircraft carriers involved in operations in the Philippines until summoned to Iwo Jima to bombard the island in advance of the Marine landings. Afterward, she proceeded to Okinawa, bombarding the island in advance of the Allied amphibious assault. In mid-1945 Wisconsin turned her attention to bombarding the Japanese home islands until the surrender of Japan in August. Reactivated in 1950, for the Korean War, Wisconsin served two tours of duty, assisting South Korean and UN forces by providing call fire support and shelling targets. In 1956, the bow of the uncompleted Kentucky was removed and grafted on Wisconsin, which had collided with the destroyer . Decommissioned in 1958, Wisconsin was placed in the reserve fleet at the Philadelphia Naval Shipyard until reactivated in 1986 as part of the 600-ship Navy plan. In 1991, Wisconsin participated in Operation Desert Storm, firing 24 Tomahawk Missiles at Iraqi targets and expending 319 16-inch shells at Iraqi troop formations along the coast. Decommissioned for the last time 30 September 1991, Wisconsin was placed in the reserve fleet until stricken from the Naval Vessel Register on 17 March 2006, so she could be transferred for use as a museum ship. Wisconsin is currently berthed at the Nauticus maritime museum in Norfolk, Virginia.
Illinois and Kentucky
Hull numbers BB-65 and BB-66 were originally intended as the first and second ships of the Montana-class of battleships; however, the passage of an emergency war building program on 19 July 1940 resulted in both hulls being reordered as Iowa-class units to save time on construction. The war ended before either could be completed, and work was eventually stopped. Initially, proposals were made to convert the hulls into aircraft carriers similar to the Essex class, but the effort was dropped.
was ordered on 9 September 1940 and initially laid down on 6 December 1942. However, work was suspended pending a decision on whether to convert the hull to an aircraft carrier. Upon determination the result would cost more and be less capable than building from scratch, construction resumed, but it was canceled for good approximately one-quarter complete on 11 August 1945. She was sold for scrap and broken up on the slipway in September 1958.
was ordered on 9 September 1940 and laid down on 7 March 1942. Work on the ship was suspended in June 1942, and the hull floated out to make room for the construction of LSTs. The interruption lasted for two and a half years while a parallel aircraft carrier debate played out as with Illinois, reaching the same conclusion. Work resumed in December 1944, with completion projected for mid-1946. Further suggestions were made to convert Kentucky into a specialist anti-aircraft ship, and work was again suspended. With the hull approximately three-quarters completed, she was floated on 20 January 1950, to clear a dry dock for repairs to Missouri, which had run aground. During this period, plans were proposed to convert Kentucky into a guided missile battleship, which saw her reclassified from BB-66 to BBG-1. When these failed construction of any sort, work never resumed and the ship was used as a parts hulk; in 1956, her bow was removed and shipped in one piece across Hampton Roads and grafted onto Wisconsin, which had collided with the destroyer Eaton. In 1958, the engines installed on Kentucky were salvaged and installed on the Sacramento-class fast combat support ships and . Ultimately, what remained of the hulk was sold for scrap on 31 October 1958.
| Technology | Naval warfare | null |
264740 | https://en.wikipedia.org/wiki/Isomerization | Isomerization | In chemistry, isomerization or isomerisation is the process in which a molecule, polyatomic ion or molecular fragment is transformed into an isomer with a different chemical structure. Enolization is an example of isomerization, as is tautomerization. When the isomerization occurs intramolecularly it may be called a rearrangement reaction.
When the activation energy for the isomerization reaction is sufficiently small, both isomers will exist in a temperature-dependent equilibrium with each other. Many values of the standard free energy difference, , have been calculated, with good agreement between observed and calculated data.
Examples and applications
Alkanes
Skeletal isomerization occurs in the cracking process, used in the petrochemical industry to convert straight chain alkanes to isoparaffins as exemplified in the conversion of normal octane to 2,5-dimethylhexane (an "isoparaffin"):
Fuels containing branched hydrocarbons are favored for internal combustion engines for their higher octane rating.
Alkenes
Trans-alkenes are about 1 kcal/mol more stable than cis-alkenes. An example of this effect is cis- vs trans-2-butene. The difference is attributed to unfavorable non-bonded interactions in the cis isomer.
Terminal alkenes isomerize to internal alkenes in the presence of metal catalysts. This process is employed in the Shell higher olefin process to convert alpha-olefins to internal olefins, which are subjected to olefin metathesis. In certain kinds of alkene polymerization reactions, chain walking is an isomerization process that introduces branches into growing polymers.
The trans isomer of resveratrol can be converted to the cis isomer in a photochemical reaction.
Thermal rearrangement of azulene to naphthalene has been observed.
Other examples
Aldose-ketose isomerism, also known as Lobry de Bruyn–van Ekenstein transformation, provides an example in saccharide chemistry.
An example of an organometallic isomerization is the production of decaphenylferrocene, from its linkage isomer.
| Physical sciences | Basics_3 | Chemistry |
265044 | https://en.wikipedia.org/wiki/Propellant | Propellant | A propellant (or propellent) is a mass that is expelled or expanded in such a way as to create a thrust or another motive force in accordance with Newton's third law of motion, and "propel" a vehicle, projectile, or fluid payload. In vehicles, the engine that expels the propellant is called a reaction engine. Although technically a propellant is the reaction mass used to create thrust, the term "propellant" is often used to describe a substance which contains both the reaction mass and the fuel that holds the energy used to accelerate the reaction mass. For example, the term "propellant" is often used in chemical rocket design to describe a combined fuel/propellant, although the propellants should not be confused with the fuel that is used by an engine to produce the energy that expels the propellant. Even though the byproducts of substances used as fuel are also often used as a reaction mass to create the thrust, such as with a chemical rocket engine, propellant and fuel are two distinct concepts.
Vehicles can use propellants to move by ejecting a propellant backwards which creates an opposite force that moves the vehicle forward. Projectiles can use propellants that are expanding gases which provide the motive force to set the projectile in motion. Aerosol cans use propellants which are fluids that are compressed so that when the propellant is allowed to escape by releasing a valve, the energy stored by the compression moves the propellant out of the can and that propellant forces the aerosol payload out along with the propellant. Compressed fluid may also be used as a simple vehicle propellant, with the potential energy that is stored in the compressed fluid used to expel the fluid as the propellant. The energy stored in the fluid was added to the system when the fluid was compressed, such as compressed air. The energy applied to the pump or thermal system that is used to compress the air is stored until it is released by allowing the propellant to escape. Compressed fluid may also be used only as energy storage along with some other substance as the propellant, such as with a water rocket, where the energy stored in the compressed air is the fuel and the water is the propellant.
In electrically powered spacecraft, electricity is used to accelerate the propellant. An electrostatic force may be used to expel positive ions, or the Lorentz force may be used to expel negative ions and electrons as the propellant. Electrothermal engines use the electromagnetic force to heat low molecular weight gases (e.g. hydrogen, helium, ammonia) into a plasma and expel the plasma as propellant. In the case of a resistojet rocket engine, the compressed propellant is simply heated using resistive heating as it is expelled to create more thrust.
In chemical rockets and aircraft, fuels are used to produce an energetic gas that can be directed through a nozzle, thereby producing thrust. In rockets, the burning of rocket fuel produces an exhaust, and the exhausted material is usually expelled as a propellant under pressure through a nozzle. The exhaust material may be a gas, liquid, plasma, or a solid. In powered aircraft without propellers such as jets, the propellant is usually the product of the burning of fuel with atmospheric oxygen so that the resulting propellant product has more mass than the fuel carried on the vehicle.
Proposed photon rockets would use the relativistic momentum of photons to create thrust. Even though photons do not have mass, they can still act as a propellant because they move at relativistic speed, i.e., the speed of light. In this case Newton's third Law of Motion is inadequate to model the physics involved and relativistic physics must be used.
In chemical rockets, chemical reactions are used to produce energy which creates movement of a fluid which is used to expel the products of that chemical reaction (and sometimes other substances) as propellants. For example, in a simple hydrogen/oxygen engine, hydrogen is burned (oxidized) to create and the energy from the chemical reaction is used to expel the water (steam) to provide thrust. Often in chemical rocket engines, a higher molecular mass substance is included in the fuel to provide more reaction mass.
Rocket propellant may be expelled through an expansion nozzle as a cold gas, that is, without energetic mixing and combustion, to provide small changes in velocity to spacecraft by the use of cold gas thrusters, usually as maneuvering thrusters.
To attain a useful density for storage, most propellants are stored as either a solid or a liquid.
Vehicle propellants
A rocket propellant is a mass that is expelled from a vehicle, such as a rocket, in such a way as to create a thrust in accordance with Newton's third law of motion, and "propel" the vehicle forward. The engine that expels the propellant is called a reaction engine. Although the term "propellant" is often used in chemical rocket design to describe a combined fuel/propellant, propellants should not be confused with the fuel that is used by an engine to produce the energy that expels the propellant. Even though the byproducts of substances used as fuel are also often used as a reaction mass to create the thrust, such as with a chemical rocket engine, propellant and fuel are two distinct concepts.
In electrically powered spacecraft, electricity is used to accelerate the propellant. An electrostatic force may be used to expel positive ions, or the Lorentz force may be used to expel negative ions and electrons as the propellant. Electrothermal engines use the electromagnetic force to heat low molecular weight gases (e.g. hydrogen, helium, ammonia) into a plasma and expel the plasma as propellant. In the case of a resistojet rocket engine, the compressed propellant is simply heated using resistive heating as it is expelled to create more thrust.
In chemical rockets and aircraft, fuels are used to produce an energetic gas that can be directed through a nozzle, thereby producing thrust. In rockets, the burning of rocket fuel produces an exhaust, and the exhausted material is usually expelled as a propellant under pressure through a nozzle. The exhaust material may be a gas, liquid, plasma, or a solid. In powered aircraft without propellers such as jets, the propellant is usually the product of the burning of fuel with atmospheric oxygen so that the resulting propellant product has more mass than the fuel carried on the vehicle.
The propellant or fuel may also simply be a compressed fluid, with the potential energy that is stored in the compressed fluid used to expel the fluid as the propellant. The energy stored in the fluid was added to the system when the fluid was compressed, such as compressed air. The energy applied to the pump or thermal system that is used to compress the air is stored until it is released by allowing the propellant to escape. Compressed fluid may also be used only as energy storage along with some other substance as the propellant, such as with a water rocket, where the energy stored in the compressed air is the fuel and the water is the propellant.
Proposed photon rockets would use the relativistic momentum of photons to create thrust. Even though photons do not have mass, they can still act as a propellant because they move at relativistic speed, i.e., the speed of light. In this case Newton's third Law of Motion is inadequate to model the physics involved and relativistic physics must be used.
In chemical rockets, chemical reactions are used to produce energy which creates movement of a fluid which is used to expel the products of that chemical reaction (and sometimes other substances) as propellants. For example, in a simple hydrogen/oxygen engine, hydrogen is burned (oxidized) to create and the energy from the chemical reaction is used to expel the water (steam) to provide thrust. Often in chemical rocket engines, a higher molecular mass substance is included in the fuel to provide more reaction mass.
Rocket propellant may be expelled through an expansion nozzle as a cold gas, that is, without energetic mixing and combustion, to provide small changes in velocity to spacecraft by the use of cold gas thrusters, usually as maneuvering thrusters.
To attain a useful density for storage, most propellants are stored as either a solid or a liquid.
Propellants may be energized by chemical reactions to expel solid, liquid or gas. Electrical energy may be used to expel gases, plasmas, ions, solids or liquids. Photons may be used to provide thrust via relativistic momentum.
Chemically powered
Solid propellant
Composite propellants made from a solid oxidizer such as ammonium perchlorate or ammonium nitrate, a synthetic rubber such as HTPB, PBAN, or Polyurethane (or energetic polymers such as polyglycidyl nitrate or polyvinyl nitrate for extra energy), optional high-explosive fuels (again, for extra energy) such as RDX or nitroglycerin, and usually a powdered metal fuel such as aluminum.
Some amateur propellants use potassium nitrate, combined with sugar, epoxy, or other fuels and binder compounds.
Potassium perchlorate has been used as an oxidizer, paired with asphalt, epoxy, and other binders.
Propellants that explode in operation are of little practical use currently, although there have been experiments with Pulse Detonation Engines. Also the newly synthesized bishomocubane based compounds are under consideration in the research stage as both solid and liquid propellants of the future.
Grain
Solid fuel/propellants are used in forms called grains. A grain is any individual particle of fuel/propellant regardless of the size or shape. The shape and size of a grain determines the burn time, amount of gas, and rate of produced energy from the burning of the fuel and, as a consequence, thrust vs time profile.
There are three types of burns that can be achieved with different grains.
Progressive burn Usually a grain with multiple perforations or a star cut in the center providing a lot of surface area.
Degressive burn Usually a solid grain in the shape of a cylinder or sphere.
Neutral burn Usually a single perforation; as outside surface decreases the inside surface increases at the same rate.
Composition
There are four different types of solid fuel/propellant compositions:
Single-based fuel/propellant A single based fuel/propellant has nitrocellulose as its chief explosives ingredient. Stabilizers and other additives are used to control the chemical stability and enhance its properties.
Double-based fuel/propellant Double-based fuel/propellants consist of nitrocellulose with nitroglycerin or other liquid organic nitrate explosives added. Stabilizers and other additives are also used. Nitroglycerin reduces smoke and increases the energy output. Double-based fuel/propellants are used in small arms, cannons, mortars and rockets.
Triple-based fuel/propellant Triple-based fuel/propellants consist of nitrocellulose, nitroguanidine, nitroglycerin or other liquid organic nitrate explosives. Triple-based fuel/propellants are used in cannons.
Composite Composites do not utilize nitrocellulose, nitroglycerin, nitroguanidine or any other organic nitrate as the primary constituent. Composites usually consist of a fuel such as metallic aluminum, a combustible binder such as synthetic rubber or HTPB, and an oxidizer such as ammonium perchlorate. Composite fuel/propellants are used in large rocket motors. In some applications, such as the US SLBM Trident II missile, nitroglycerin is added to the aluminum and ammonium perchlorate composite as an energetic plasticizer.
Liquid propellant
In rockets, three main liquid bipropellant combinations are used: cryogenic oxygen and hydrogen, cryogenic oxygen and a hydrocarbon, and storable propellants.
Cryogenic oxygen-hydrogen combination system Used in upper stages and sometimes in booster stages of space launch systems. This is a nontoxic combination. This gives high specific impulse and is ideal for high-velocity missions.
Cryogenic oxygen-hydrocarbon propellant system Used for many booster stages of space launch vehicles as well as a smaller number of second stages. This combination of fuel/oxidizer has high density and hence allows for a more compact booster design.
Storable propellant combinations Used in almost all bipropellant low-thrust, auxiliary or reaction control rocket engines, as well as in some in large rocket engines for first and second stages of ballistic missiles. They are instant-starting and suitable for long-term storage.
Propellant combinations used for liquid propellant rockets include:
Liquid oxygen and liquid hydrogen
Liquid oxygen and kerosene or RP-1
Liquid oxygen and ethanol
Liquid oxygen and methane
Hydrogen peroxide and mentioned above alcohol or RP-1
Red fuming nitric acid (RFNA) and kerosene or RP-1
RFNA and Unsymmetrical dimethylhydrazine (UDMH)
Dinitrogen tetroxide and UDMH, MMH, and/or hydrazine
Common monopropellant used for liquid rocket engines include:
Hydrogen peroxide
Hydrazine
Red fuming nitric acid (RFNA)
Electrically powered
Electrically powered reactive engines use a variety of usually ionized propellants, including atomic ions, plasma, electrons, or small droplets or solid particles as propellant.
Electrostatic
If the acceleration is caused mainly by the Coulomb force (i.e. application of a static electric field in the direction of the acceleration) the device is considered electrostatic. The types of electrostatic drives and their propellants:
Gridded ion thruster – using positive ions as the propellant, accelerated by an electrically charged grid
NASA Solar Technology Application Readiness (NSTAR) – positive ions accelerated using high-voltage electrodes
HiPEP – using positive ions as the propellant, created using microwaves
Radiofrequency ion thruster – generalization of HiPEP
Hall-effect thruster, including its subtypes Stationary Plasma Thruster (SPT) and Thruster with Anode Layer (TAL) – use the Hall effect to orient electrons to create positive ions for propellant
Colloid ion thruster – electrostatic acceleration of droplets of liquid salt as the propellant
Field-emission electric propulsion – using electrodes to accelerate ionized liquid metal as a propellant
Nano-particle field extraction thruster – using charged cylindrical carbon nanotubes as propellant
Electrothermal
These are engines that use electromagnetic fields to generate a plasma which is used as the propellant. They use a nozzle to direct the energized propellant. The nozzle itself may be composed simply of a magnetic field. Low molecular weight gases (e.g. hydrogen, helium, ammonia) are preferred propellants for this kind of system.
Resistojet – using a usually inert compressed propellant that is energized by simple resistive heating
Arcjet – uses (usually) hydrazine or ammonia as a propellant which is energized with an electrical arc
Microwave – a type of Radiofrequency ion thruster
Variable specific impulse magnetoplasma rocket (VASIMR) – using microwave-generated plasma as the propellant and magnetic field to direct its expulsion
Electromagnetic
Electromagnetic thrusters use ions as the propellant, which are accelerated by the Lorentz force or by magnetic fields, either of which is generated by electricity:
Electrodeless plasma thruster – a complex system that uses cold plasma as a propellant that is accelerated by ponderomotive force
Magnetoplasmadynamic thruster – propellants include xenon, neon, argon, hydrogen, hydrazine, or lithium; expelled using the Lorentz force
Pulsed inductive thruster – because this reactive engine uses a radial magnetic field, it acts on both positive and negative particles and so it may use a wide range of gases as a propellant including water, hydrazine, ammonia, argon, xenon and many others
Pulsed plasma thruster – uses a Teflon plasma as a propellant, which is created by an electrical arc and expelled using the Lorentz force
Helicon Double Layer Thruster – a plasma propellant is generated and excited from a gas using a helicon induced by high frequency band radiowaves which form a magnetic nozzle in a cylinder
Nuclear
Nuclear reactions may be used to produce the energy for the expulsion of the propellants. Many types of nuclear reactors have been used/proposed to produce electricity for electrical propulsion as outlined above. Nuclear pulse propulsion uses a series of nuclear explosions to create large amounts of energy to expel the products of the nuclear reaction as the propellant. Nuclear thermal rockets use the heat of a nuclear reaction to heat a propellant. Usually the propellant is hydrogen because the force is a function of the energy irrespective of the mass of the propellant, so the lightest propellant (hydrogen) produces the greatest specific impulse.
Photonic
A photonic reactive engine uses photons as the propellant and their discrete relativistic energy to produce thrust.
Projectile propellants
Compressed fluid propellants
Compressed fluid or compressed gas propellants are pressurized physically, by a compressor, rather than by a chemical reaction. The pressures and energy densities that can be achieved, while insufficient for high-performance rocketry and firearms, are adequate for most applications, in which case compressed fluids offer a simpler, safer, and more practical source of propellant pressure.
A compressed fluid propellant may simply be a pressurized gas, or a substance which is a gas at atmospheric pressure, but stored under pressure as a liquid.
Compressed gas propellants
In applications in which a large quantity of propellant is used, such as pressure washing and airbrushing, air may be pressurized by a compressor and used immediately. Additionally, a hand pump to compress air can be used for its simplicity in low-tech applications such as atomizers, plant misters and water rockets. The simplest examples of such a system are squeeze bottles for such liquids as ketchup and shampoo.
However, compressed gases are impractical as stored propellants if they do not liquify inside the storage container, because very high pressures are required in order to store any significant quantity of gas, and high-pressure gas cylinders and pressure regulators are expensive and heavy.
Liquified gas propellants
Principle
Liquefied gas propellants are gases at atmospheric pressure, but become liquid at a modest pressure. This pressure is high enough to provide useful propulsion of the payload (e.g. aerosol paint, deodorant, lubricant), but is low enough to be stored in an inexpensive metal can, and to not pose a safety hazard in case the can is ruptured.
The mixture of liquid and gaseous propellant inside the can maintains a constant pressure, called the liquid's vapor pressure. As the payload is depleted, the propellant vaporizes to fill the internal volume of the can. Liquids are typically 500-1000x denser than their corresponding gases at atmospheric pressure; even at the higher pressure inside the can, only a small fraction of its volume needs to be propellant in order to eject the payload and replace it with vapor.
Vaporizing the liquid propellant to gas requires some energy, the enthalpy of vaporization, which cools the system. This is usually insignificant, although it can sometimes be an unwanted effect of heavy usage (as the system cools, the vapor pressure of the propellant drops). However, in the case of a freeze spray, this cooling contributes to the desired effect (although freeze sprays may also contain other components, such as chloroethane, with a lower vapor pressure but higher enthalpy of vaporization than the propellant).
Propellant compounds
Chlorofluorocarbons (CFCs) were once often used as propellants, but since the Montreal Protocol came into force in 1989, they have been replaced in nearly every country due to the negative effects CFCs have on Earth's ozone layer. The most common replacements of CFCs are mixtures of volatile hydrocarbons, typically propane, n-butane and isobutane. Dimethyl ether (DME) and methyl ethyl ether are also used. All these have the disadvantage of being flammable. Nitrous oxide and carbon dioxide are also used as propellants to deliver foodstuffs (for example, whipped cream and cooking spray). Medicinal aerosols such as asthma inhalers use hydrofluoroalkanes (HFA): either HFA 134a (1,1,1,2,-tetrafluoroethane) or HFA 227 (1,1,1,2,3,3,3-heptafluoropropane) or combinations of the two. More recently, liquid hydrofluoroolefin (HFO) propellants have become more widely adopted in aerosol systems due to their relatively low vapor pressure, low global warming potential (GWP), and nonflammability.
Payloads
The practicality of liquified gas propellants allows for a broad variety of payloads. Aerosol sprays, in which a liquid is ejected as a spray, include paints, lubricants, degreasers, and protective coatings; deodorants and other personal care products; cooking oils. Some liquid payloads are not sprayed due to lower propellant pressure and/or viscous payload, as with whipped cream and shaving cream or shaving gel. Low-power guns, such as BB guns, paintball guns, and airsoft guns, have solid projectile payloads. Uniquely, in the case of a gas duster ("canned air"), the only payload is the velocity of the propellant vapor itself.
| Technology | Energy: General | null |
265062 | https://en.wikipedia.org/wiki/Salt%20marsh | Salt marsh | A salt marsh, saltmarsh or salting, also known as a coastal salt marsh or a tidal marsh, is a coastal ecosystem in the upper coastal intertidal zone between land and open saltwater or brackish water that is regularly flooded by the tides. It is dominated by dense stands of salt-tolerant plants such as herbs, grasses, or low shrubs. These plants are terrestrial in origin and are essential to the stability of the salt marsh in trapping and binding sediments. Salt marshes play a large role in the aquatic food web and the delivery of nutrients to coastal waters. They also support terrestrial animals and provide coastal protection.
Salt marshes have historically been endangered by poorly implemented coastal management practices, with land reclaimed for human uses or polluted by upstream agriculture or other industrial coastal uses. Additionally, sea level rise caused by climate change is endangering other marshes, through erosion and submersion of otherwise tidal marshes. However, recent acknowledgment by both environmentalists and larger society for the importance of saltwater marshes for biodiversity, ecological productivity and other ecosystem services, such as carbon sequestration, have led to an increase in salt marsh restoration and management since the 1980s.
Basic information
Salt marshes occur on low-energy shorelines in temperate and high-latitudes which can be stable, emerging, or submerging depending if the sedimentation is greater, equal to, or lower than relative sea level rise (subsidence rate plus sea level change), respectively. Commonly these shorelines consist of mud or sand flats (known also as tidal flats or abbreviated to mudflats) which are nourished with sediment from inflowing rivers and streams. These typically include sheltered environments such as embankments, estuaries and the leeward side of barrier islands and spits. In the tropics and sub-tropics they are replaced by mangroves; an area that differs from a salt marsh in that instead of herbaceous plants, they are dominated by salt-tolerant trees.
Most salt marshes have a low topography with low elevations but a vast wide area, making them hugely popular for human populations.
Salt marshes are located among different landforms based on their physical and geomorphological settings. Such marsh landforms include deltaic marshes, estuarine, back-barrier, open coast, embayments and drowned-valley marshes. Deltaic marshes are associated with large rivers where many occur in Southern Europe such as the Camargue, France in the Rhône delta or the Ebro delta in Spain. They are also extensive within the rivers of the Mississippi River Delta in the United States. In New Zealand, most salt marshes occur at the head of estuaries in areas where there is little wave action and high sedimentation. Such marshes are located in Awhitu Regional Park in Auckland, the Manawatū Estuary, and the Avon Heathcote Estuary / Ihutai in Christchurch. Back-barrier marshes are sensitive to the reshaping of barriers in the landward side of which they have been formed. They are common along much of the eastern coast of the United States and the Frisian Islands. Large, shallow coastal embayments can hold salt marshes with examples including Morecambe Bay and Portsmouth in Britain and the Bay of Fundy in North America.
Salt marshes are sometimes included in lagoons, and the difference is not very marked; the Venetian Lagoon in Italy, for example, is made up of these sorts of animals and or living organisms belonging to this ecosystem. They have a big impact on the biodiversity of the area. Salt marsh ecology involves complex food webs which include primary producers (vascular plants, macroalgae, diatoms, epiphytes, and phytoplankton), primary consumers (zooplankton, macrozoa, molluscs, insects), and secondary consumers.
The low physical energy and high grasses provide a refuge for animals. Many marine fish use salt marshes as nursery grounds for their young before they move to open waters. Birds may raise their young among the high grasses, because the marsh provides both sanctuary from predators and abundant food sources which include fish trapped in pools, insects, shellfish, and worms.
Worldwide occurrence
Saltmarshes across 99 countries (essentially worldwide) were mapped by Mcowen et al. 2017. A total of 5,495,089 hectares of mapped saltmarsh across 43 countries and territories are represented in a Geographic Information Systems polygon shapefile. This estimate is at the relatively low end of previous estimates (2.2–40 Mha). A later study conservatively estimated global saltmarsh extent as 90,800 km2 (9,080,000 hectares). The most extensive saltmarshes worldwide are found outside the tropics, notably including the low-lying, ice-free coasts, bays and estuaries of the North Atlantic which are well represented in their global polygon dataset.
Formation
The formation begins as tidal flats gain elevation relative to sea level by sediment accretion, and subsequently the rate and duration of tidal flooding decreases so that vegetation can colonize on the exposed surface. The arrival of propagules of pioneer species such as seeds or rhizome portions are combined with the development of suitable conditions for their germination and establishment in the process of colonisation. When rivers and streams arrive at the low gradient of the tidal flats, the discharge rate reduces and suspended sediment settles onto the tidal flat surface, helped by the backwater effect of the rising tide. Mats of filamentous blue-green algae can fix silt and clay sized sediment particles to their sticky sheaths on contact which can also increase the erosion resistance of the sediments. This assists the process of sediment accretion to allow colonising species (e.g., Salicornia spp.) to grow. These species retain sediment washed in from the rising tide around their stems and leaves and form low muddy mounds which eventually coalesce to form depositional terraces, whose upward growth is aided by a sub-surface root network which binds the sediment. Once vegetation is established on depositional terraces further sediment trapping and accretion can allow rapid upward growth of the marsh surface such that there is an associated rapid decrease in the depth and duration of tidal flooding. As a result, competitive species that prefer higher elevations relative to sea level can inhabit the area and often a succession of plant communities develops.
Tidal flooding and vegetation zonation
Coastal salt marshes can be distinguished from terrestrial habitats by the daily tidal flow that occurs and continuously floods the area. It is an important process in delivering sediments, nutrients and plant water supply to the marsh. At higher elevations in the upper marsh zone, there is much less tidal inflow, resulting in lower salinity levels. Soil salinity in the lower marsh zone is fairly constant due to everyday annual tidal flow. However, in the upper marsh, variability in salinity is shown as a result of less frequent flooding and climate variations. Rainfall can reduce salinity and evapotranspiration can increase levels during dry periods. As a result, there are microhabitats populated by different species of flora and fauna dependent on their physiological abilities. The flora of a salt marsh is differentiated into levels according to the plants' individual tolerance of salinity and water table levels. Vegetation found at the water must be able to survive high salt concentrations, periodical submersion, and a certain amount of water movement, while plants further inland in the marsh can sometimes experience dry, low-nutrient conditions. It has been found that the upper marsh zones limit species through competition and the lack of habitat protection, while lower marsh zones are determined through the ability of plants to tolerate physiological stresses such as salinity, water submergence and low oxygen levels.
The New England salt marsh is subject to strong tidal influences and shows distinct patterns of zonation. In low marsh areas with high tidal flooding, a monoculture of the smooth cordgrass, Spartina alterniflora dominate, then heading landwards, zones of the salt hay, Spartina patens, black rush, Juncus gerardii and the shrub Iva frutescens are seen respectively. These species all have different tolerances that make the different zones along the marsh best suited for each individual.
Plant species diversity is relatively low, since the flora must be tolerant of salt, complete or partial submersion, and anoxic mud substrate. The most common salt marsh plants are glassworts (Salicornia spp.) and the cordgrass (Spartina spp.), which have worldwide distribution. They are often the first plants to take hold in a mudflat and begin its ecological succession into a salt marsh. Their shoots lift the main flow of the tide above the mud surface while their roots spread into the substrate and stabilize the sticky mud and carry oxygen into it so that other plants can establish themselves as well. Plants such as sea lavenders (Limonium spp.), plantains (Plantago spp.), and varied sedges and rushes grow once the mud has been vegetated by the pioneer species.
Salt marshes are quite photosynthetically active and are extremely productive habitats. They serve as depositories for a large amount of organic matter and are full of decomposition, which feeds a broad food chain of organisms from bacteria to mammals. Many of the halophytic plants such as cordgrass are not grazed at all by higher animals but die off and decompose to become food for micro-organisms, which in turn become food for fish and birds.
Sediment trapping, accretion, and the role of tidal creeks
The factors and processes that influence the rate and spatial distribution of sediment accretion within the salt marsh are numerous. Sediment deposition can occur when marsh species provide a surface for the sediment to adhere to, followed by deposition onto the marsh surface when the sediment flakes off at low tide. The amount of sediment adhering to salt marsh species is dependent on the type of marsh species, the proximity of the species to the sediment supply, the amount of plant biomass, and the elevation of the species. For example, in a study of the Eastern Chongming Island and Jiuduansha Island tidal marshes at the mouth of the Yangtze River, China, the amount of sediment adhering to the species Spartina alterniflora, Phragmites australis, and Scirpus mariqueter decreased with distance from the highest levels of suspended sediment concentrations (found at the marsh edge bordering tidal creeks or the mudflats); decreased with those species at the highest elevations, which experienced the lowest frequency and depth of tidal inundations; and increased with increasing plant biomass. Spartina alterniflora, which had the most sediment adhering to it, may contribute >10% of the total marsh surface sediment accretion by this process.
Salt marsh species also facilitate sediment accretion by decreasing current velocities and encouraging sediment to settle out of suspension. Current velocities can be reduced as the stems of tall marsh species induce hydraulic drag, with the effect of minimising re-suspension of sediment and encouraging deposition. Measured concentrations of suspended sediment in the water column have been shown to decrease from the open water or tidal creeks adjacent to the marsh edge, to the marsh interior, probably as a result of direct settling to the marsh surface by the influence of the marsh canopy.
Inundation and sediment deposition on the marsh surface is also assisted by tidal creeks which are a common feature of salt marshes. Their typically dendritic and meandering forms provide avenues for the tide to rise and flood the marsh surface, as well as to drain water, and they may facilitate higher amounts of sediment deposition than salt marsh bordering open ocean. Sediment deposition is correlated with sediment size: coarser sediments will deposit at higher elevations (closer to the creek) than finer sediments (further from the creek). Sediment size is also often correlated with particular trace metals, and thus tidal creeks can affect metal distributions and concentrations in salt marshes, in turn affecting the biota. Salt marshes do not however require tidal creeks to facilitate sediment flux over their surface although salt marshes with this morphology seem to be rarely studied.
The elevation of marsh species is important; those species at lower elevations experience longer and more frequent tidal floods and therefore have the opportunity for more sediment deposition to occur. Species at higher elevations can benefit from a greater chance of inundation at the highest tides when increased water depths and marsh surface flows can penetrate into the marsh interior.
Human impacts
The coast is a highly attractive natural feature to humans through its beauty, resources, and accessibility. As of 2002, over half of the world's population was estimated to being living within 60 km of the coastal shoreline, making coastlines highly vulnerable to human impacts from daily activities that put pressure on these surrounding natural environments. In the past, salt marshes were perceived as coastal 'wastelands,' causing considerable loss and change of these ecosystems through land reclamation for agriculture, urban development, salt production and recreation. The indirect effects of human activities such as nitrogen loading also play a major role in the salt marsh area. Salt marshes can suffer from dieback in the high marsh and die-off in the low marsh. A study published in 2022 estimates that 22% of saltmarsh loss from 1999–2019 was due to direct human drivers, defined as observable activities occurring at the location of the detected change, such as conversion to aquaculture, agriculture, coastal development, or other physical structures. Additionally, 30% of saltmarsh gain over this same time period were also due to direct drivers, such as restoration activities or coastal modifications to promote tidal exchange.
Land reclamation
Reclamation of land for agriculture by converting marshland to upland was historically a common practice. Dikes were often built to allow for this shift in land change and to provide flood protection further inland. In recent times intertidal flats have also been reclaimed. For centuries, livestock such as sheep and cattle grazed on the highly fertile salt marsh land. Land reclamation for agriculture has resulted in many changes such as shifts in vegetation structure, sedimentation, salinity, water flow, biodiversity loss and high nutrient inputs. There have been many attempts made to eradicate these problems for example, in New Zealand, the cordgrass Spartina anglica was introduced from England into the Manawatū River mouth in 1913 to try and reclaim the estuary land for farming. A shift in structure from bare tidal flat to pastureland resulted from increased sedimentation and the cordgrass extended out into other estuaries around New Zealand. Native plants and animals struggled to survive as non-natives out competed them. Efforts are now being made to remove these cordgrass species, as the damages are slowly being recognized.
In the Blyth estuary in Suffolk in eastern England, the mid-estuary reclamations (Angel and Bulcamp marshes) that were abandoned in the 1940s have been replaced by tidal flats with compacted soils from agricultural use overlain with a thin veneer of mud. Little vegetation colonisation has occurred in the last 60–75 years and has been attributed to a combination of surface elevations too low for pioneer species to develop, and poor drainage from the compacted agricultural soils acting as an aquiclude. Terrestrial soils of this nature need to adjust from fresh to saline interstitial water by a change in the chemistry and the structure of the soil, accompanied with fresh deposition of estuarine sediment, before salt marsh vegetation can establish. The vegetation structure, species richness, and plant community composition of salt marshes naturally regenerated on reclaimed agricultural land can be compared to adjacent reference salt marshes to assess the success of marsh regeneration.
Upstream agriculture
Cultivation of land upstream from the salt marsh can introduce increased silt inputs and raise the rate of primary sediment accretion on the tidal flats, so that pioneer species can spread further onto the flats and grow rapidly upwards out of the level of tidal inundation. As a result, marsh surfaces in this regime may have an extensive cliff at their seaward edge. At the Plum Island estuary, Massachusetts (U.S.), stratigraphic cores revealed that during the 18th and 19th century the marsh prograded over subtidal and mudflat environments to increase in area from 6 km2 to 9 km2 after European settlers deforested the land upstream and increased the rate of sediment supply.
Urban development and nitrogen loading
The conversion of marshland to upland for agriculture has in the past century been overshadowed by conversion for urban development. Coastal cities worldwide have encroached onto former salt marshes and in the U.S. the growth of cities looked to salt marshes for waste disposal sites. Estuarine pollution from organic, inorganic, and toxic substances from urban development or industrialisation is a worldwide problem and the sediment in salt marshes may entrain this pollution with toxic effects on floral and faunal species. Urban development of salt marshes has slowed since about 1970 owing to growing awareness by environmental groups that they provide beneficial ecosystem services. They are highly productive ecosystems, and when net productivity is measured in g m−2 yr−1 they are equalled only by tropical rainforests. Additionally, they can help reduce wave erosion on sea walls designed to protect low-lying areas of land from wave erosion.
De-naturalisation of the landward boundaries of salt marshes from urban or industrial encroachment can have negative effects. In the Avon-Heathcote estuary/Ihutai, New Zealand, species abundance and the physical properties of the surrounding margins were strongly linked, and the majority of salt marsh was found to be living along areas with natural margins in the Avon / Ōtākaro and Ōpāwaho / Heathcote river outlets; conversely, artificial margins contained little marsh vegetation and restricted landward retreat. The remaining marshes surrounding these urban areas are also under immense pressure from the human population as human-induced nitrogen enrichment enters these habitats. Nitrogen loading through human-use indirectly affects salt marshes causing shifts in vegetation structure and the invasion of non-native species.
Human impacts such as sewage, urban run-off, agricultural and industrial wastes are running into the marshes from nearby sources. Salt marshes are nitrogen limited and with an increasing level of nutrients entering the system from anthropogenic effects, the plant species associated with salt marshes are being restructured through change in competition. For example, the New England salt marsh is experiencing a shift in vegetation structure where S. alterniflora is spreading from the lower marsh where it predominately resides up into the upper marsh zone. Additionally, in the same marshes, the reed Phragmites australis has been invading the area expanding to lower marshes and becoming a dominant species. P. australis is an aggressive halophyte that can invade disturbed areas in large numbers outcompeting native plants. This loss in biodiversity is not only seen in flora assemblages but also in many animals such as insects and birds as their habitat and food resources are altered.
Sea level rise
Due to the melting of Arctic sea ice and thermal expansion of the oceans, as a result of global warming, sea levels have begun to rise. As with all coastlines, this rise in water levels is predicted to negatively affect salt marshes, by flooding and eroding them. The sea level rise causes more open water zones within the salt marsh. These zones cause erosion along their edges, further eroding the marsh into open water until the whole marsh disintegrates.
While salt marshes are susceptible to threats concerning sea level rise, they are also an extremely dynamic coastal ecosystem. Salt marshes may in fact have the capability to keep pace with a rising sea level, by 2100, mean sea level could see increases between 0.6m to 1.1m. Marshes are susceptible to both erosion and accretion, which play a role in a what is called a bio-geomorphic feedback. Salt marsh vegetation captures sediment to stay in the system which in turn allows for the plants to grow better and thus the plants are better at trapping sediment and accumulate more organic matter. This positive feedback loop potentially allows for salt marsh bed level rates to keep pace with rising sea level rates. However, this feedback is also dependent on other factors like productivity of the vegetation, sediment supply, land subsidence, biomass accumulation, and magnitude and frequency of storms. In a study published by Ü. S. N. Best in 2018, they found that bioaccumulation was the number one factor in a salt marsh's ability to keep up with SLR rates. The salt marsh's resilience depends upon its increase in bed level rate being greater than that of sea levels' increasing rate, otherwise the marsh will be overtaken and drowned.
Biomass accumulation can be measured in the form of above-ground organic biomass accumulation, and below-ground inorganic accumulation by means of sediment trapping and sediment settling from suspension. Salt marsh vegetation helps to increase sediment settling because it slows current velocities, disrupts turbulent eddies, and helps to dissipate wave energy. Marsh plant species are known for their tolerance to increased salt exposure due to the common inundation of marshlands. These types of plants are called halophytes. Halophytes are a crucial part of salt marsh biodiversity and their potential to adjust to elevated sea levels. With elevated sea levels, salt marsh vegetation would likely be more exposed to more frequent inundation rates and it must be adaptable or tolerant to the consequential increased salinity levels and anaerobic conditions. There is a common elevation (above the sea level) limit for these plants to survive, where anywhere below the optimal line would lead to anoxic soils due to constant submergence and too high above this line would mean harmful soil salinity levels due to the high rate of evapotranspiration as a result of decreased submergence.
Along with the vertical accretion of sediment and biomass, the accommodation space for marsh land growth must also be considered. Accommodation space is the land available for additional sediments to accumulate and marsh vegetation to colonize laterally. This lateral accommodation space is often limited by anthropogenic structures such as coastal roads, sea walls and other forms of development of coastal lands. A study by Lisa M. Schile, published in 2014, found that across a range of sea level rise rates, marshlands with high plant productivity were resistant against sea level rises but all reached a pinnacle point where accommodation space was necessary for continued survival. The presence of accommodation space allows for new mid/high habitats to form, and for marshes to escape complete inundation.
Mosquito control
Earlier in the 20th century, it was believed that draining salt marshes would help reduce mosquito populations, such as Aedes taeniorhynchus, the black salt marsh mosquito. In many locations, particularly in the northeastern United States, residents and local and state agencies dug straight-lined ditches deep into the marsh flats. The end result, however, was a depletion of killifish habitat. The killifish is a mosquito predator, so the loss of habitat actually led to higher mosquito populations, and adversely affected wading birds that preyed on the killifish. These ditches can still be seen, despite some efforts to refill the ditches.
Crab herbivory and bioturbation
Increased nitrogen uptake by marsh species into their leaves can prompt greater rates of length-specific leaf growth, and increase the herbivory rates of crabs. The burrowing crab Neohelice granulata frequents SW Atlantic salt marshes where high density populations can be found among populations of the marsh species Spartina densiflora and Sarcocornia perennis. In Mar Chiquita lagoon, north of Mar del Plata, Argentina, Neohelice granulata herbivory increased as a likely response to the increased nutrient value of the leaves of fertilised Spartina densiflora plots, compared to non-fertilised plots. Regardless of whether the plots were fertilised or not, grazing by Neohelice granulata also reduced the length specific leaf growth rates of the leaves in summer, while increasing their length-specific senescence rates. This may have been assisted by the increased fungal effectiveness on the wounds left by the crabs.
The salt marshes of Cape Cod, Massachusetts (US), are experiencing creek bank die-offs of Spartina spp. (cordgrass) that has been attributed to herbivory by the crab Sesarma reticulatum. At 12 surveyed Cape Cod salt marsh sites, 10% – 90% of creek banks experienced die-off of cordgrass in association with a highly denuded substrate and high density of crab burrows. Populations of Sesarma reticulatum are increasing, possibly as a result of the degradation of the coastal food web in the region. The bare areas left by the intense grazing of cordgrass by Sesarma reticulatum at Cape Cod are suitable for occupation by another burrowing crab, Uca pugnax, which are not known to consume live macrophytes. The intense bioturbation of salt marsh sediments from this crab's burrowing activity has been shown to dramatically reduce the success of Spartina alterniflora and Suaeda maritima seed germination and established seedling survival, either by burial or exposure of seeds, or uprooting or burial of established seedlings. However, bioturbation by crabs may also have a positive effect. In New Zealand, the tunnelling mud crab Helice crassa has been given the stately name of an 'ecosystem engineer' for its ability to construct new habitats and alter the access of nutrients to other species. Their burrows provide an avenue for the transport of dissolved oxygen in the burrow water through the oxic sediment of the burrow walls and into the surrounding anoxic sediment, which creates the perfect habitat for special nitrogen cycling bacteria. These nitrate reducing (denitrifying) bacteria quickly consume the dissolved oxygen entering into the burrow walls to create the oxic mud layer that is thinner than that at the mud surface. This allows a more direct diffusion path for the export of nitrogen (in the form of gaseous nitrogen (N2)) into the flushing tidal water.
Microbial life in salt marshes
The variable salinity, climate, nutrient levels and anaerobic conditions of salt marshes provide strong selective pressures on the microorganisms inhabiting them. In salt marshes, microbes play the main role in nutrient cycling and biogeochemical processing. To date, the microbial community of salt marshes has not been found to change drastically due to human impacts, but the research is still ongoing. Because of the major role of microbes in these environments, it is critical to understand the different processes performed and different microbial players present in salt marshes. Salt marshes provide habitat for chemo(litho)autotrophs, heterotrophs, and photoautotrophs alike. These organisms contribute diverse environmental services such as sulfate reduction, nitrification, decomposition and rhizosphere interactions.
Chemo(litho)autotrophs in salt marshes
Chemoautotrophs, also known as chemolithoautotrophs, are organisms capable of creating their own energy, from the use of inorganic molecules, and are able to thrive in harsh environments, such as deep sea vents or salt marshes, due to not depending upon external organic carbon sources for their growth and survival. Some Chemoautotrophic bacterial microorganisms found in salt marshes include Betaproteobacteria and Gammaproteobacteria, both classes including sulfate-reducing bacteria (SRB), sulfur-oxidizing bacteria (SOB), and ammonia-oxidizing bacteria (AOB) which play crucial roles in nutrient cycling and ecosystem functioning.
Abundance and diversity of sulfate-reducing chemolithoautotrophs
Bacterial chemolithoautotrophs in salt marshes include sulfate-reducing bacteria. In these ecosystems, up to 50% of sedimentary remineralization can be attributed to sulfate reduction. The dominant class of sulfate-reducing bacteria in salt marshes tends to be Deltaproteobacteria. Some examples of deltaproteobacteria that are found in salt marshes are species of genera Desulfobulbus, Desulfuromonas, and Desulfovibrio.
The abundance and diversity of chemolithoautotrophs in salt marshes is largely determined by the composition of plant species in the salt marsh ecosystem. Each type of salt-marsh plant has varying lengths of growing seasons, varying photosynthetic rates, and they all lose varying amounts of organic matter to the ocean, resulting in varying carbon-inputs to the ecosystem. The results from an experiment that was done in a salt marsh in the Yangtze estuary in China, suggested that both the species richness and total abundance of sulfate-reducing bacterial communities increased when a new plant, S. alterniflora, with a higher C-input to the ecosystem was introduced. Although chemolithotrophs produce their own carbon, they still depend on the C-input from salt marshes because of the indirect impact it has on the amount of viable electron donors, such as reduced sulfur compounds. The concentration of reduced sulfur compounds, as well as other possible electron donors, increases with more organic-matter decomposition (by other organisms). Therefore if the ecosystem contains more decomposing organic matter, as with plants with high photosynthetic and littering rates, there will be more electron donors available to the bacteria, and thus more sulfate reduction is possible. As a result, the abundance of sulfate-reducing bacteria increases. The high-photosynthetic-rate, high-litter-rate salt marsh plant, S. alterniflora, was discovered to withstand high sulfur concentrations in the soil, which would normally be somewhat toxic to plants.
The abundance of chemolithoautotrophs in salt marshes also varies temporally as a result of being somewhat dependent on the organic C-input from plants in the ecosystem. Since plants grow most throughout the summer, and usually begin to lose biomass around fall during their late stage, the highest input of decomposing organic matter is in the fall. Thus seasonally, the abundance of chemolithotrophs in salt marshes is highest in autumn.
Why are sulfate-reducing bacteria in salt marshes?
Salt marshes are the ideal environment for sulfate-reducing bacteria. The sulfate-reducing bacteria tend to live in anoxic conditions, such as in salt marshes, because they require reduced compounds to produce their energy. Since there is a high sedimentation rate and a high amount of organic matter, the conditions of the sediment are usually dependably anoxic. However, the conditions all across the salt marsh (above the sediment) are not completely anoxic, which means the organisms living here must have some level of tolerance to oxygen. Many of the chemolithoautotrophs living outside or at the surface of the sediment also exhibit this characteristic.
Significance of sulfate-reducing bacteria
Sulfate-reducing bacteria play a significant role in nutrient recycling and in reducing nitrate pollution levels. Since humans have been adding disproportionate amounts of nitrates to coastal waters, salt marshes are one of the ecosystems where nitrate pollution remains an issue. The enrichment of nitrates in the water increases denitrification, as well as microbial decomposition and primary productivity. Sulfate-reducing and oxidizing bacteria, however, play a role in removing the excess nitrates from the water to prevent eutrophication. Since the sulfate-reducing bacteria is in the water and sediment, reduced sulfur molecules are usually in abundance. These reduced sulfates then react with excess nitrate in the water, reducing nitrate and oxidizing the reduced sulfur. As a result of human nitrate enrichment, it is predicted that sulfur-oxidizing bacteria which also reduce nitrates will increase in relative abundance to sulfur-reducing bacteria.
Abundance and significance of chemolithoautotroph nitrifiers within salt marshes
Within salt marshes, chemolithoautotrophic nitrifying bacteria are also frequently identified, including Betaproteobacteria ammonia oxidizers such as Nitrosomonas and Nitrosospira. Although ammonia-oxidizing Archaea (AOA) are found to be more prevalent than ammonium-oxidizing Bacteria (AOB) within salt marsh environments, predominantly from the Crenarchaeota group, AOB play a critical role within the salt marsh environment too. Increases in marsh salinity tend to favor AOB, while higher oxygen levels and lower carbon-to-nitrogen ratios favor AOA. These AOB are important in catalyzing the rate-limiting step within the nitrification process, by using ammonium monooxygenase (AMO), produced from amoA, to convert ammonium (NH4+) into nitrite (NO2-). Specifically, within the class of Betaproteobacteria, Nitrosomonas aestuarii, Nitrosomonas marina, and Nitrosospira ureae are highly prevalent within the salt marsh environment; similarly, within the class of Gammaproteobacteria, Nitrosococcus spp. are key AOB in the marshes.
The abundance of these chemolithoautotrophs varies along the salinity gradients present within salt marshes: Nitrosomonas are more prevalent within lower salinity or freshwater regions, while Nitrosospira are found to dominate in higher saline environments. In addition, the abundance of fixed-nitrogen in these environments critically influences the distribution of the betaproteobacteria within the salt marsh: Nitrosomonas are more found to be in greater abundance within high N and C environments, whereas Nitrosospira are found to be more abundant in lower N and C regions. Further, factors such as temperature, pH, net primary productivity, and regions of anoxia may limit nitrification, and thus critically influence nitrifier distribution.
The role of nitrification by AOB in salt marshes critically links ammonia, produced from the mineralization of organic nitrogen compounds, to the process of nitrogen oxidation. Further, nitrogen oxidation is important for the downstream removal of nitrates into nitrogen gas, catalyzed by denitrifiers, from the marsh environment. Hence, AOB play an indirect role in nitrogen removal into the atmosphere.
Photoautotrophic bacteria
The bacterial photoautotroph community of salt marshes primarily consists of cyanobacteria, purple bacteria, and green sulfur bacteria.
Cyanobacteria in salt marshes
Cyanobacteria are important nitrogen fixers in salt marshes, and provide nitrogen to organisms like diatoms and microalgae.
Purple bacteria
Oxygen inhibits photosynthesis in purple bacteria, which makes estuaries a favorable habitat for them due to the low oxygen content and high levels of light present, optimizing their photosynthesis. In anoxic environments, like salt marshes, many microbes have to use sulfate as an electron acceptor during cellular respiration instead of oxygen, producing lots of hydrogen sulfide as a byproduct. While hydrogen sulfide is toxic to most organisms, purple bacteria require it to grow and will metabolize it to either sulfate or sulfur, and by doing so allowing other organisms to inhabit the toxic environment. Purple bacteria can be further classified as either purple sulphur bacteria, or purple non-sulfur bacteria. Purple sulphur bacteria are more tolerant to sulfide and store the sulfur they create intracellularly, while purple non-sulfur bacteria excrete any sulfur they produce.
Green bacteria
Green sulfur bacteria (Chlorobiaceae) are photoautotrophic bacteria that utilize sulfide and thiosulfate for their growth, producing sulfate in the process. They are very adapted to photosynthesizing in low light environments with bacteriochlorophyll pigments a, c, d, and e, to help them absorb wavelengths of light that other organisms cannot. When co-existing with purple bacteria, they often occupy lower depths as they are less tolerant to oxygen, but more photosynthetically adept.
Rhizosphere microbes
Fungi
Some mycorrhizal fungi, like arbuscular mycorrhiza are widely associated with salt marsh plants and may even help plants grow in salt marsh soil rich in heavy metals by reducing their uptake into the plant, although the exact mechanism has yet to be determined.
Bacteria
Examining 16S ribosomal DNA found in Yangtze River Estuary, the most common bacteria in the rhizosphere were Proteobacteria such as Betaproteobacteria, Gammaproteobacteria, Deltaproteobacteria, and Epsilonproteobacteria. One such widespread species had a similar ribotype to the animal pathogen S. marcescens, and may be beneficial for plants as the bacteria can break down chitin into available carbon and nitrogen for plants to use. Actinobacteria have also been found in plant rhizosphere in costal salt marshes and help plants grow through helping plants absorb more nutrients and secreting antimicrobial compounds. In Jiangsu, China, Actinobacteria from the suborders Pseudonocardineae, Corynebacterineae, Propionibacterineae, Streptomycineae, Micromonosporineae, Streptosporangineae and Micrococcineae were cultured and isolated from rhizosphere soil.
Microbial decomposition activity within salt marshes
Another key process among microbial salt marshes is microbial decomposition activity. Nutrient cycling in salt marshes is highly promoted by the resident community of bacteria and fungi involved in remineralizing organic matter. Studies on the decomposition of a salt marsh cordgrass, Spartina alterniflora, have shown that fungal colonization begins the degradation process, which is then finished by the bacterial community. The carbon from Spartina alterniflora is made accessible to the salt marsh food web largely through these bacterial communities which are then consumed by bacterivores. Bacteria are responsible for the degradation of up to 88% of lignocellulotic material in salt marshes. However, fungal populations have been found to dominate over bacterial populations in winter months.
The fungi that make up the decomposition community in salt marshes come from the phylum ascomycota, the two most prevalent species being Phaeosphaeria spartinicola and Mycosphaerella sp. strain 2. In terms of bacteria, the alphaproteobacteria class is the most prevalent class within the salt marsh environment involved in decomposition activity. The propagation of Phaeosphaeria spartinicola is through ascospores that are released when the host plant is wetted by high tides or rain.
Restoration and management
The perception of bay salt marshes as a coastal 'wasteland' has since changed, acknowledging that they are one of the most biologically productive habitats on earth, rivalling tropical rainforests. Salt marshes are ecologically important, providing habitats for native migratory fish and acting as sheltered feeding and nursery grounds. They are now protected by legislation in many countries to prevent the loss of these ecologically important habitats. In the United States and Europe, they are now accorded a high level of protection by the Clean Water Act and the Habitats Directive respectively. With the impacts of this habitats and their importance now realised, a growing interest in restoring salt marshes through managed retreat or the reclamation of land has been established. However, many Asian countries such as China still need to recognise the value of marshlands. With their ever-growing populations and intense development along the coast, the value of salt marshes tends to be ignored and the land continues to be reclaimed.
Bakker et al. (1997) suggests two options available for restoring salt marshes. The first is to abandon all human interference and leave the salt marsh to complete its natural development. These types of restoration projects are often unsuccessful as vegetation tends to struggle to revert to its original structure and the natural tidal cycles are shifted due to land changes. The second option suggested by Bakker et al. (1997) is to restore the destroyed habitat into its natural state either at the original site or as a replacement at a different site. Under natural conditions, recovery can take 2–10 years or even longer depending on the nature and degree of the disturbance and the relative maturity of the marsh involved. Marshes in their pioneer stages of development will recover more rapidly than mature marshes as they are often first to colonize the land. It is important to note that restoration can often be sped up through the replanting of native vegetation.
This last approach is often the most practiced and generally more successful than allowing the area to naturally recover on its own. The salt marshes in the state of Connecticut in the United States have long been an area lost to fill and dredging. As of 1969, the Tidal Wetland Act was introduced that ceased this practice, but despite the introduction of the act, the system was still degrading due to alterations in tidal flow. One area in Connecticut is the marshes on Barn Island. These marshes were diked then impounded with salt and brackish marsh during 1946–1966. As a result, the marsh shifted to a freshwater state and became dominated by the invasive species P. australis, Typha angustifolia and T. latifolia that have little ecological connection to the area.
By 1980, a restoration programme was put in place that has now been running for over 20 years. This programme has aimed to reconnect the marshes by returning tidal flow along with the ecological functions and characteristics of the marshes back to their original state. In the case of Barn Island, reduction of the invasive species has been initiated, re-establishing the tidal-marsh vegetation along with animal species such as fish and insects. This example highlights that considerable time and effort is needed to effectively restore salt marsh systems. The timescale for salt marsh recovery is dependent on the development stage of the marsh, type and extent of the disturbance, geographical location and the environmental and physiological stress factors to the marsh-associated flora and fauna.
Although much effort has gone into restoring salt marshes worldwide, further research is needed. There are many setbacks and problems associated with marsh restoration that require careful long-term monitoring. Information on all components of the salt marsh ecosystem should be understood and monitored from sedimentation, nutrient, and tidal influences, to behaviour patterns and tolerances of both flora and fauna species. Once a better understanding of these processes is acquired, and not just locally, but over a global scale, then more sound and practical management and restoration efforts can be implemented to preserve these valuable marshes and restore them to their original state.
While humans are situated along coastlines, there will always be the possibility of human-induced disturbances despite the number of restoration efforts we plan to implement. Dredging, pipelines for offshore petroleum resources, highway construction, accidental toxic spills or just plain carelessness are examples that will for some time now and into the future be the major influences of salt marsh degradation.
In addition to restoring and managing salt marsh systems based on scientific principles, the opportunity should be taken to educate public audiences of their importance biologically and their purpose as serving as a natural buffer for flood protection. Because salt marshes are often located next to urban areas, they are likely to receive more visitors than remote wetlands. By physically seeing the marsh, people are more likely to take notice and be more aware of the environment around them. An example of public involvement occurred at the Famosa Slough State Marine Conservation Area in San Diego, where a "friends" group worked for over a decade in trying to prevent the area from being developed. Eventually, the site was bought by the city and the group worked together to restore the area. The project involved removing of invasive species and replanting with native ones, along with public talks to other locals, frequent bird walks and clean-up events.
Research methods
There is a diverse range and combination of methodologies employed to understand the hydrological dynamics in salt marshes and their ability to trap and accrete sediment. Sediment traps are often used to measure rates of marsh surface accretion when short term deployments (e.g. less than one month) are required. These circular traps consist of pre-weighed filters that are anchored to the marsh surface, then dried in a laboratory and re-weighed to determine the total deposited sediment.
For longer term studies (e.g. more than one year) researchers may prefer to measure sediment accretion with marker horizon plots. Marker horizons consist of a mineral such as feldspar that is buried at a known depth within wetland substrates to record the increase in overlying substrate over long time periods. In order to gauge the amount of sediment suspended in the water column, manual or automated samples of tidal water can be poured through pre-weighed filters in a laboratory then dried to determine the amount of sediment per volume of water.
Another method for estimating suspended sediment concentrations is by measuring the turbidity of the water using optical backscatter probes, which can be calibrated against water samples containing a known suspended sediment concentration to establish a regression relationship between the two. Marsh surface elevations may be measured with a stadia rod and transit, electronic theodolite, Real-Time Kinematic Global Positioning System, laser level or electronic distance meter (total station). Hydrological dynamics include water depth, measured automatically with a pressure transducer, or with a marked wooden stake, and water velocity, often using electromagnetic current meters.
| Physical sciences | Wetlands | Earth science |
265112 | https://en.wikipedia.org/wiki/Improvised%20explosive%20device | Improvised explosive device | An improvised explosive device (IED) is a bomb constructed and deployed in ways other than in conventional military action. It may be constructed of conventional military explosives, such as an artillery shell, attached to a detonating mechanism. IEDs are commonly used as roadside bombs, or homemade bombs.
The term "IED" was coined by the British Army during the Northern Ireland conflict to refer to booby traps made by the IRA, and entered common use in the U.S. during the Iraq War.
IEDs are generally utilized in terrorist operations or in asymmetric unconventional warfare or urban warfare by insurgent guerrillas or commando forces in a theatre of operations. In the Iraq War (2003–2011), insurgents used IEDs extensively against U.S.-led forces, and by the end of 2007, IEDs were responsible for approximately 63% of coalition deaths in Iraq. They were also used in Afghanistan by insurgent groups, and caused over 66% of coalition casualties in the 2001–2021 Afghanistan War.
IEDs were also used frequently by the Liberation Tigers of Tamil Eelam (LTTE) in Sri Lanka during the Sri Lankan Civil War, by the Chechen insurgency following the Second Chechen War, and by Ambazonian separatists in the ongoing Anglophone Crisis.
Background
An IED is a bomb fabricated in an improvised manner incorporating destructive, lethal, noxious, pyrotechnic, or incendiary chemicals and designed to destroy or incapacitate personnel or vehicles. In some cases, IEDs are used to distract, disrupt, or delay an opposing force, facilitating another type of attack. IEDs may incorporate military or commercially sourced explosives, and often combine both types, or they may otherwise be made with homemade explosives (HME). An HME lab refers to a Homemade Explosive Lab, or the physical location where the devices are crafted.
An IED has five components: a switch (activator), an initiator (fuse), container (body), charge (explosive), and a power source (battery). An IED designed for use against armoured targets such as personnel carriers or tanks will be designed for armour penetration, by using a shaped charge that creates an explosively formed penetrator. IEDs are extremely diverse in design and may contain many types of initiators, detonators, penetrators, and explosive loads.
Antipersonnel IEDs typically also contain fragmentation-generating objects such as nails, ball bearings or even small rocks to cause wounds at greater distances than blast pressure alone could. In the conflicts of the 21st century, anti-personnel improvised explosive devices (IED) have partially replaced conventional or military landmines as the source of injury to dismounted (pedestrian) soldiers and civilians. These injuries were reported in BMJ Open to be far worse with IEDs than with landmines resulting in multiple limb amputations and lower body mutilation. This combination of injuries has been given the name "Dismounted Complex Blast Injury" and is thought to be the worst survivable injury ever seen in war.
IEDs are triggered by various methods, including remote control, infrared or magnetic triggers, pressure-sensitive bars or trip wires (victim-operated). In some cases, multiple IEDs are wired together in a daisy chain to attack a convoy of vehicles spread out along a roadway.
IEDs made by inexperienced designers or with substandard materials may fail to detonate, and in some cases, they detonate on either the maker or the placer of the device. Some groups, however, have been known to produce sophisticated devices constructed with components scavenged from conventional munitions and standard consumer electronics components, such as mobile phones, consumer-grade two-way radios, washing machine timers, pagers, or garage door openers. The sophistication of an IED depends on the training of the designer and the tools and materials available.
IEDs may use artillery shells or conventional high-explosive charges as their explosive load as well as homemade explosives. However, the threat exists that toxic chemical, biological, or radioactive (dirty bomb) material may be added to a device, thereby creating other life-threatening effects beyond the shrapnel, concussive blasts and fire normally associated with bombs. Chlorine liquid has been added to IEDs in Iraq, producing clouds of chlorine gas.
A vehicle-borne IED, or VBIED, is a military term for a car bomb or truck bomb but can be any type of transportation such as a bicycle, motorcycle, donkey (), etc. They are typically employed by insurgents, in particular ISIS, and can carry a relatively large payload. They can also be detonated from a remote location. VBIEDs can create additional shrapnel through the destruction of the vehicle itself and use vehicle fuel as an incendiary weapon. The act of a person's being in this vehicle and detonating it is known as an SVBIED suicide.
Of increasing popularity among insurgent forces in Iraq is the house-borne IED, or HBIED, from the common military practice of clearing houses; insurgents rig an entire house to detonate and collapse shortly after a clearing squad has entered.
By warhead
The Dictionary of Military and Associated Terms (JCS Pub 1-02) includes two definitions for improvised devices: improvised explosive devices (IED) and improvised nuclear device (IND). These definitions address the Nuclear and Explosive in CBRNe. That leaves chemical, biological and radiological undefined. Four definitions have been created to build on the structure of the JCS definition. Terms have been created to standardize the language of first responders and members of the military and to correlate the operational picture.
Explosive
A device placed or fabricated in an improvised manner incorporating destructive, lethal, noxious, pyrotechnic, or incendiary chemicals and designed to destroy, incapacitate, harass, or distract. It may incorporate military stores, but is normally devised from non-military components.
Explosively formed penetrator/projectiles (EFPs)
IEDs have been deployed in the form of explosively formed projectiles (EFP), a special type of shaped charge that is effective at long standoffs from the target (50 meters or more), however they are not accurate at long distances. This is because of how they are produced. The large "slug" projected from the explosion has no stabilization because it has no tail fins and it does not spin like a bullet from a rifle. Without this stabilization the trajectory can not be accurately determined beyond 50 meters. An EFP is essentially a cylindrical shaped charge with a machined concave metal disc (often copper) in front, pointed inward. The force of the shaped charge turns the disc into a high velocity slug, capable of penetrating the armor of most vehicles in Iraq.
Directionally focused charges
Directionally focused charges (also known as directionally focused fragmentary charges depending on the construction) are very similar to EFPs, with the main difference being that the top plate is usually flat and not concave. It also is not made with machined copper but much cheaper cast or cut metal. When made for fragmentation, the contents of the charge are usually nuts, bolts, ball bearings and other similar shrapnel products and explosive. If it only consists of the flat metal plate, it is known as a platter charge, serving a similar role as an EFP with reduced effect but easier construction.
Chemical
A device incorporating the toxic attributes of chemical materials designed to result in the dispersal of toxic chemical materials for the purpose of creating a primary patho-physiological toxic effect (morbidity and mortality), or secondary psychological effect (causing fear and behavior modification) on a larger population. Such devices may be fabricated in a completely improvised manner or may be an improvised modification to an existing weapon.
Biological
A device incorporating biological materials designed to result in the dispersal of vector borne biological material for the purpose of creating a primary patho-physiological toxic effect (morbidity and mortality), or secondary psychological effect (causing fear and behavior modification) on a larger population.
Incendiary
A device making use of exothermic chemical reactions designed to result in the rapid spread of fire for the purpose of creating a primary patho-physiological effect (morbidity and mortality), or secondary psychological effect (causing fear and behavior modification) on a larger population or it may be used with the intent of gaining a tactical advantage. Such devices may be fabricated in a completely improvised manner or may be an improvised modification to an existing weapon. A common type of this is the Molotov cocktail.
Radiological
A speculative device incorporating radioactive materials designed to result in the dispersal of radioactive material for the purpose of area denial and economic damage, and/or for the purpose of creating a primary patho-physiological toxic effect (morbidity and mortality), or secondary psychological effect (causing fear and behavior modification) on a larger population. Such devices may be fabricated in a completely improvised manner or may be an improvised modification to an existing nuclear weapon. Also called a Radiological Dispersion Device (RDD) or "dirty bomb".
Nuclear
Improvised nuclear device of most likely gun-type or implosion-type.
By delivery mechanism
Car
A vehicle may be laden with explosives, set to explode by remote control or by a passenger/driver, commonly known as a car bomb or vehicle-borne IED (VBIED, pronounced vee-bid). On occasion the driver of the car bomb may have been coerced into delivery of the vehicle under duress, a situation known as a proxy bomb. Distinguishing features are low-riding vehicles with excessive weight, vehicles with only one passenger, and ones where the interior of the vehicles look as if they have been stripped down and built back up. Car bombs can carry thousands of pounds of explosives and may be augmented with shrapnel to increase fragmentation.
ISIS has used truck bombs with devastating effects.
Boat (WBIED)
Water-borne Improvised Explosive Devices (WBIED), i.e. boats carrying explosives, can be used against ships and areas connected to water.
An early example of this type was the Japanese Shinyo suicide boats during World War II. The boats were filled with explosives and attempted to ram Allied ships, sometimes successfully, having sunk or severely damaged several American ships by war's end. Suicide bombers used a boat-borne IED to attack the USS Cole; US and UK troops have also been killed by boat-borne IEDs in Iraq. The Tamil Tigers Sea Tigers have also been known to use SWBIEDs during the Sri Lankan Civil War.
WBIEDs have been used in the Red Sea.
Animal
Monkeys and war pigs were used as incendiaries around 1000 AD. More famously the "anti-tank dog" and "bat bomb" were developed during World War II. In recent times, a two-year-old child and seven other people were killed by explosives strapped to a horse in the town of Chita in Colombia. The carcasses of certain animals were also used to conceal explosive devices by the Iraqi insurgency.
Collar
IEDs strapped to the necks of farmers have been used on at least three occasions by guerrillas in Colombia, as a way of extortion. American pizza delivery man Brian Douglas Wells was killed in 2003 by an explosive fastened to his neck, purportedly under duress from the maker of the bomb. In 2011 a schoolgirl in Sydney, Australia had a suspected collar bomb attached to her by an attacker in her home. The device was removed by police after a ten-hour operation and proved to be a hoax.
Suicide
Suicide bombing usually refers to an individual wearing explosives and detonating them to kill others including themselves, the bomber will conceal explosives on and around their person, commonly using a vest, and will use a timer or some other trigger to detonate the explosives. The logic behind such attacks is the belief that an IED delivered by a human has a greater chance of achieving success than any other method of attack. In addition, there is the psychological impact of child soldiers prepared to deliberately sacrifice themselves for their cause.
Surgically implanted
In May 2012 American counter-terrorism officials leaked their acquisition of documents describing the preparation and use of surgically implanted improvised explosive devices.
The devices were designed to evade detection.
The devices were described as containing no metal, so they could not be detected by X-rays.
Security officials referred to bombs being surgically implanted into suicide bombers' "love handles".
According to the Daily Mirror UK security officials at MI-6 asserted that female bombers could travel undetected carrying the explosive chemicals in otherwise standard breast implants. The bomber would blow up the implanted explosives by injecting a chemical trigger.
Robot
Robots could also be used to carry explosives. First such documented case was during the aftermath of 2016 shooting of Dallas police officers when a bomb disposal robot was used to deliver explosives to kill Micah Xavier Johnson, who was hiding in a place inaccessible to police snipers. As well, drones carrying explosives were used in a suspected assassination attempt against Venezuelan president Nicolás Maduro in 2018.
Tunnel
ISIS
and Al-Nusra have used bombs detonated in tunnels dug under targets.
Improvised rocket
In 2008, rocket-propelled IEDs, dubbed Improvised Rocket Assisted Munitions, Improvised Rocket Assisted Mortars and (IRAM) by the military, came to be employed in numbers against U.S. forces in Iraq. They have been described as propane tanks packed with explosives and powered by 107 mm rockets. They are similar to some Provisional IRA barrack buster mortars.
New types of IRAMs including Volcano IRAM and Elephant Rockets, are used during the Syrian Civil War.
Improvised mortar
Improvised mortars have been used by many insurgent groups including during the civil war in Syria and Boko Haram insurgency. IRA used improvised mortars called barrack busters.
Improvised artillery including hell cannons are used by rebel forces during Syrian Civil War.
By trigger mechanism
Wire
Command-wire improvised, explosive devices (CWIED) use an electrical firing cable that affords the user complete control over the device right up until the moment of initiation.
Radio
The trigger for a radio-controlled improvised explosive device (RCIED) is controlled by radio link. The device is constructed so that the receiver is connected to an electrical firing circuit and the transmitter operated by the perpetrator at a distance. A signal from the transmitter causes the receiver to trigger a firing pulse that operates the switch. Usually the switch fires an initiator; however, the output may also be used to remotely arm an explosive circuit. Often the transmitter and receiver operate on a matched coding system that prevents the RCIED from being initiated by spurious radio frequency signals or jamming. An RCIED can be triggered from any number of different radio-frequency based mechanisms including handheld remote control transmitters, car alarms, wireless door bells, cell phones, pagers and portable two-way radios, including those designed for the UHF PMR446, FRS, and GMRS services.
Mobile phone
A radio-controlled IED (RCIED) incorporating a mobile phone that is modified and connected to an electrical firing circuit. Mobile phones operate in the UHF band in line of sight with base transceiver station (BTS) antennae sites. In the common scenario, receipt of a paging signal by phone is sufficient to initiate the IED firing circuit.
Victim-operated
Victim-operated improvised explosive devices (VOIED), also known as booby traps, are designed to function upon contact with a victim. VOIED switches are often well hidden from the victim or disguised as innocuous everyday objects. They are operated by means of movement. Switching methods include tripwire, pressure mats, spring-loaded release, push, pull or tilt. Common forms of VOIED include the under-vehicle IED (UVIED), improvised landmines, and mail bombs.
Infrared
The British accused Iran and Hezbollah of teaching Iraqi fighters to use infrared light beams to trigger IEDs. As the occupation forces became more sophisticated in interrupting radio signals around their convoys, the insurgents adapted their triggering methods. In some cases, when a more advanced method was disrupted, the insurgents regressed to using uninterruptible means, such as hard wires from the IED to detonator; however, this method is much harder to effectively conceal. It later emerged however, that these "advanced" IEDs were actually old IRA technology. The infrared beam method was perfected by the IRA in the early 1990s after it acquired the technology from a botched undercover British Army operation. Many of the IEDs being used against the invading coalition forces in Iraq were originally developed by the British Army who unintentionally passed the information on to the IRA. The IRA taught their techniques to the Palestine Liberation Organisation and the knowledge spread to Iraq.
Counterefforts
Counter-IED efforts are done primarily by military, law enforcement, diplomatic, financial, and intelligence communities and involve a comprehensive approach to countering the threat networks that employ IEDs, not just efforts to defeat the devices themselves.
Detection and disarmament
Because the components of these devices are being used in a manner not intended by their manufacturer, and because the method of producing the explosion is limited only by the science and imagination of the perpetrator, it is not possible to follow a step-by-step guide to detect and disarm a device that an individual has only recently developed. As such, explosive ordnance disposal (IEDD) operators must be able to fall back on their extensive knowledge of the first principles of explosives and ammunition, to try and deduce what the perpetrator has done, and only then to render it safe and dispose of or exploit the device.
Beyond this, as the stakes increase and IEDs are emplaced not only to achieve the direct effect, but to deliberately target IEDD operators and cordon personnel, the IEDD operator needs to have a deep understanding of tactics to ensure they are neither setting up any of their team or the cordon troops for an attack, nor walking into one themselves. The presence of chemical, biological, radiological, or nuclear (CBRN) material in an IED requires additional precautions. As with other missions, the EOD operator provides the area commander with an assessment of the situation and of support needed to complete the mission.
Military and law enforcement personnel from around the world have developed a number of render-safe procedures (RSPs) to deal with IEDs. RSPs may be developed as a result of direct experience with devices or by applied research designed to counter the threat. The supposed effectiveness of IED jamming systems, including vehicle- and personally-mounted systems, has caused IED technology to essentially regress to command-wire detonation methods. These are physical connections between the detonator and explosive device and cannot be jammed. However, these types of IEDs are more difficult to emplace quickly, and are more readily detected.
Military forces and law enforcement from India, Canada, United Kingdom, Israel, Spain, and the United States are at the forefront of counter-IED efforts, as all have direct experience in dealing with IEDs used against them in conflict or terrorist attacks. From the research and development side, programs such as the new Canadian Unmanned Systems Challenge will bring student groups together to invent an unmanned device to both locate IEDs and pinpoint the insurgents.
Historical use
The fougasse was improvised for centuries, eventually inspiring factory-made land mines. Ernst Jünger mentions in his war memoir the systematic use of IEDs and booby traps to cover the retreat of German troops at the Somme region during World War I. Another early example of coordinated large-scale use of IEDs was the Belarusian Rail War launched by Belarusian guerrillas against the Germans during World War II. Both command-detonated and delayed-fuse IEDs were used to derail thousands of German trains during 1943–1944.
Afghanistan
Starting six months before the invasion of Afghanistan by the USSR on 27 December 1979, the Afghan Mujahideen were supplied by the CIA, among others, with large quantities of military supplies. Among those supplies were many types of anti-tank mines. The insurgents often removed the explosives from several foreign anti-tank mines, and combined the explosives in tin cooking-oil cans for a more powerful blast. By combining the explosives from several mines and placing them in tin cans, the insurgents made them more powerful, but sometimes also easier to detect by Soviet sappers using mine detectors. After an IED was detonated, the insurgents often used direct-fire weapons such as machine guns and rocket-propelled grenades to continue the attack.
Afghan insurgents operating far from the border with Pakistan did not have a ready supply of foreign anti-tank mines. They preferred to make IEDs from Soviet unexploded ordnance. The devices were rarely triggered by pressure fuses. They were almost always remotely detonated. Since the 2001 invasion of Afghanistan, the Taliban and its supporters have used IEDs against NATO and Afghan military and civilian vehicles. This has become the most common method of attack against NATO forces, with IED attacks increasing consistently year on year.
A brigade commander said that sniffer dogs are the most reliable way of detecting IEDs. However, statistical evidence gathered by the US Army Maneuver Support Center at Fort Leonard Wood, MO shows that the dogs are not the most effective means of detecting IEDs. The U.S. Army's 10th Mountain Division was the first unit to introduce explosive detection dogs in southern Afghanistan. In less than two years the dogs discovered 15 tons of illegal munitions, IED's, and weapons.
In July 2012 it was reported that "sticky bombs", magnetically adhesive IED's that were prevalent in the Iraq War, showed up in Afghanistan. By 2021 there was at least one sticky bomb attack a day in Kabul. They are used in both traditional assassinations and targeted killings and as terror weapons against the population at large.
In November 2013 one of the largest IEDs constructed was intercepted near Gardez City in Eastern Afghanistan. The 61,000 pounds of explosives was hidden under what appeared to be piles of wood. By comparison, the truck bomb that all but razed the Alfred P. Murrah Federal Building in Oklahoma City and killed 168 people in 1995 weighed less than 5,000 pounds. A United States Army Corps of Engineers officer assigned to the nearby FOB Lightning analyzed the potential blast damage, which resulted in closing FOB Goode due to its proximity to the highway.
ISAF troops stationed in Afghanistan and other IED prone areas of operation would commonly "BIP" (blow in place) IED's and other explosives that were considered too dangerous to defuse.
Egypt
IEDs are being used by insurgents against government forces during the insurgency in Egypt (2013–present) and the Sinai insurgency.
India
IEDs are increasingly being used by Maoists in India.
On 13 July 2011, three IEDs were used by the Insurgency in Jammu and Kashmir to carry out a coordinated attack on the city of Mumbai, killing 19 people and injuring 130 more.
On 21 February 2013, two IEDs were used to carry out bombings in the Indian city of Hyderabad. The bombs exploded in Dilsukhnagar, a crowded shopping area of the city, within 150 metres of each other.
On 17 April 2013, two kilos of explosives used in Bangalore bomb blast at Malleshwaram area, leaving 16 injured and no fatalities. Intelligence sources have said the bomb was an Improvised Explosive Device or IED.
On 21 May 2014, Indinthakarai village supporters of the Kudankulam Nuclear Power Plant were targeted by opponents using over half a dozen crude "country-made bombs". It was further reported that there had been at least four similar bombings in Tamil Nadu during the preceding year.
On 28 December 2014, a minor explosion took place near the Coconut Grove restaurant at Church Street in Bangalore on Sunday around 8:30 pm. One woman was killed and another injured in the blast.
During the 2016 Pathankot attack, several casualties came from IEDs.
On 14 February 2019 in 2019 Pulwama attack, several casualties were reported due to IED blast.
On 29 October 2023, a series of IED explosions were used to kill 2 attendees at a Jehovah's Witnesses Convention in Kalamassery, India.
Iraq
In the 2003–2011 Iraq War, IEDs have been used extensively against Coalition forces and by the end of 2007 they have been responsible for at least 64% of Coalition deaths in Iraq.
Since the detonation of the first IED in Iraq in 2003, more than 81,000 IED attacks have occurred in the country, killing and wounding 21,200 Americans.
Beginning in July 2003, the Iraqi insurgency used IEDs to target invading coalition vehicles. According to The Washington Post, 64% of U.S. deaths in Iraq occurred due to IEDs. A French study showed that in Iraq, from March 2003 to November 2006, on a global deaths in the US-led invading coalition soldiers, were caused by IEDs, i.e. 41%. That is to say more than in the "normal fights" (1027 dead, 34%). Insurgents now use the bombs to target not only invading coalition vehicles but Iraqi police as well.
Common locations for placing these bombs on the ground include animal carcasses, soft drink cans, and boxes. Typically, they explode underneath or to the side of the vehicle to cause the maximum amount of damage. However, as vehicle armour was improved on military vehicles, insurgents began placing IEDs in elevated positions such as on road signs, utility poles, or trees, to hit less protected areas.
IEDs in Iraq may be made with artillery or mortar shells or with varying amounts of bulk or homemade explosives. Early during the Iraq war, the bulk explosives were often obtained from stored munitions bunkers to include stripping landmines of their explosives.
Despite the increased armor, IEDs are killing military personnel and civilians with greater frequency. May 2007 was the deadliest month for IED attacks thus far, with a reported 89 of the 129 invading coalition casualties coming from an IED attack. According to the Pentagon, 250,000 tons (out of 650,000 tons total) of Iraqi heavy ordnance were looted, providing a large supply of ammunition for the insurgents.
In October 2005, the UK government charged that Iran was supplying insurgents with the technological know-how to make shaped charge IEDs. Both Iranian and Iraqi government officials denied the allegations.
During the Iraqi Civil War (2014–2017), ISIL has made extensive use of suicide VBIEDs, often driven by children, elderly and disabled.
On August 27, 2023, Israeli security forces successfully foiled an attempt to smuggle Iranian-made explosives into Israel from Jordan. The thwarted smuggling operation in the Jordan Valley aimed to supply terror groups in the West Bank with explosives. Counter-smuggling efforts along the border have led to increased seizures of weapons and explosive devices.
Ireland and the United Kingdom
From 1912-1913, the Suffragettes utilised IEDs in the Suffragette bombing and arson campaign.
Throughout the Troubles, the Provisional Irish Republican Army made extensive use of IEDs in their 1969–97 campaign, much of which were made in the Republic of Ireland. They used barrack buster mortars and remote-controlled IEDs. Members of the IRA developed and counter-developed devices and tactics. IRA bombs became highly sophisticated, featuring anti-handling devices such as a mercury tilt switch or microswitches. These devices would detonate the bomb if it was moved in any way. Typically, the safety-arming device used was a clockwork Memopark timer, which armed the bomb up to 60 minutes after it was placed by completing an electrical circuit supplying power to the anti-handling device. Depending on the particular design (e.g., boobytrapped briefcase or car bomb) an independent electrical circuit supplied power to a conventional timer set for the intended time delay, e.g. 40 minutes. However, some electronic delays developed by IRA technicians could be set to accurately detonate a bomb weeks after it was hidden, which is what happened in the Brighton hotel bomb attack of 1984. Initially, bombs were detonated either by timer or by simple command wire. Later, bombs could be detonated by radio control. Initially, simple servos from radio-controlled aircraft were used to close the electrical circuit and supply power to the detonator. After the British developed jammers, IRA technicians introduced devices that required a sequence of pulsed radio codes to arm and detonate them. These were harder to jam.
The IRA as well as Ulster loyalist paramilitaries have also utilized less sophisticated devices, such as homemade grenades crudely thrown at the target. These are sometimes called "blast bombs".
Roadside bombs were extensively used by the IRA. Typically, a roadside bomb was placed in a drain or culvert along a rural road and detonated by remote control when British security forces vehicles were passing, as with the case of the 1979 Warrenpoint ambush. As a result of the use of these bombs, the British military stopped transport by road in areas such as South Armagh, and used helicopter transport instead to avoid the danger.
Most IEDs used commercial or homemade explosives made in the Republic of Ireland, with ingredients such as gelignite and ANFO either stolen in construction sites or provided for by supporters in the South, although the use of Semtex-H smuggled in from Libya in the 1980s was also common from the mid-1980s onward. Bomb Disposal teams from 321 EOD manned by Ammunition Technicians were deployed in those areas to deal with the IED threat. The IRA also used secondary devices to catch British reinforcements sent in after an initial blast as occurred in the Warrenpoint Ambush. Between 1970 and 2005, the IRA detonated 19,000 IEDs in the Northern Ireland and Britain, an average of one every 17 hours for three and a half decades, arguably making it "the biggest terrorist bombing campaign in history".
In the early 1970s, at the height of the IRA campaign, the British Army unit tasked with rendering safe IEDs, 321 EOD, sustained significant casualties while engaged in bomb disposal operations. This mortality rate was far higher than other high risk occupations such as deep sea diving, and a careful review was made of how men were selected for EOD operations. The review recommended bringing in psychometric testing of soldiers to ensure those chosen had the correct mental preparation for high risk bomb disposal duties.
The IRA came up with ever more sophisticated designs and deployments of IEDs. Booby Trap or Victim Operated IEDs (VOIEDs), became commonplace. The IRA engaged in an ongoing battle to gain the upper hand in electronic warfare with remote controlled devices. The rapid changes in development led 321 EOD to employ specialists from DERA (now Dstl, an agency of the MOD), the Royal Signals, and Military Intelligence. This approach by the British army to fighting the IRA in Northern Ireland led to the development and use of most of the modern weapons, equipment and techniques now used by EOD Operators throughout the rest of the world today.
The bomb disposal operations were led by Ammunition Technicians and Ammunition Technical Officers from 321 EOD, and were trained at the Felix Centre at the Army School of Ammunition.
Israel
IEDs have been used in many attacks by Palestinian militants and continue to be used in recent attacks.
Lebanon
The Lebanese National Resistance Front, the Popular Front for the Liberation of Palestine, other resistance groups in Lebanon, and later Hezbollah, made extensive use of IEDs to resist Israeli forces after Israel's invasion of Lebanon in 1982. Israel withdrew from Beirut, Northern Lebanon, and Mount Lebanon in 1985, whilst maintaining its occupation of Southern Lebanon. Hezbollah frequently used IEDs to attack Israeli military forces in this area up until the Israeli withdrawal, and the end of the invasion of Lebanon in May 2000.
One such bomb killed Israeli Brigadier General Erez Gerstein on 28 February 1999, the highest-ranking Israeli to die in Lebanon since Yekutiel Adam's death in 1982.
Also in the 2006 War in Lebanon, a Merkava Mark II tank was hit by a pre-positioned Hezbollah IED, killing all 4 IDF servicemen on board, the first of two IEDs to damage a Merkava tank.
Libya
Homemade IEDs are used extensively during the post-civil war violence in Libya, mostly in the city of Benghazi against police stations, cars or foreign embassies.
Nepal
IEDs were also widely used in the 10-years long civil war of the Maoists in Nepal, ranging from those bought from illicit groups in India and China, to self-made devices. Typically used devices were pressure cooker bombs, socket bombs, pipe bombs, bucket bombs, etc. The devices were used more for the act of terrorizing the urban population rather than for fatal causes, placed in front of governmental offices, street corners or road sides. Mainly, the home-made IEDs were responsible for destruction of majority of structures targeted by the Maoists and assisted greatly in spreading terror among the public.
Nigeria
Boko Haram are using IEDs during their insurgency.
Pakistan
Taliban and other insurgent groups use IEDs against police, military, security forces, and civilian targets.
Russia
IEDs have also been popular in Chechnya, where Russian forces were engaged in fighting with rebel elements. While no concrete statistics are available on this matter, bombs have accounted for many Russian deaths in both the First Chechen War (1994–1996) and the Second (1999–2009).
Somalia
Al Shabaab is using IEDs during the Somali Civil War.
Syria
During the Syrian Civil War, militant insurgents were using IEDs to attack buses, cars, trucks, tanks and military convoys. Additionally, the Syrian Air Force has used barrel bombs to attack targets in cities and other areas. Such barrel bombs consist of barrels filled with high explosives, oil, and shrapnel, and are dropped from helicopters.
Along with mines and IEDs, ISIL also used VBIEDs in Syria, including during 2017 Aleppo suicide car bombing. | Technology | Explosive weapons | null |
265118 | https://en.wikipedia.org/wiki/Elephantidae | Elephantidae | Elephantidae is a family of large, herbivorous proboscidean mammals collectively called elephants and mammoths. In some cases, all members of the family can be referred to as elephants. They are large terrestrial mammals with a snout modified into a trunk and teeth modified into tusks. Most genera and species in the family are extinct. Only two genera, Loxodonta (African elephants) and Elephas (Asian elephants), are living.
The family was first described by John Edward Gray in 1821, and later assigned to taxonomic ranks within the order Proboscidea. Elephantidae has been revised by various authors to include or exclude other extinct proboscidean genera.
Description
Elephantids are distinguished from more primitive proboscideans like gomphotheres by their teeth, which have parallel lophs, formed from the merger of the cusps found in the teeth of more primitive proboscideans, which are bound by cementum. In later elephantids, these lophs became narrow lamellae, with the number of lophs/lamellae per tooth, as well as the tooth crown height (hypsodonty) increasing over time. Elephantids chew using a proal jaw movement involving a forward stroke of the lower jaws, different from the oblique movement using side to side motion of the jaws in more primitive proboscideans. The most primitive elephantid Stegotetrabelodon had a long lower jaw with lower tusks and retained permanent premolars similar to many gomphotheres, while modern elephantids lack permanent premolars, with the lower jaw being shortened (brevirostrine) and lower tusks being absent. Elephantids are typically sexually dimorphic, with substantially larger males, with an accelerated growth rate over a longer period of time than females. Elephantidae contains some of the largest known proboscideans, with fully-grown males of some species of mammoths and Palaeoloxodon having average body masses of and respectively, making them among the largest terrestrial mammals ever. One species of Palaeoloxodon, Palaeoloxodon namadicus, has been suggested to have been possibly the largest land mammal of all time, though this remains speculative due to the fragmentary nature of known remains. In comparison to more primitive elephantimorphs like gomphotheres, the bodies of elephantids tend to be proportionally shorter from front to back, as well having more elongate limbs with more slender limb bones.
Classification
Some authors have suggested to classify the family into two subfamilies, Stegotetrabelodontinae, which is monotypic, only containing Stegotetrabelodon, and Elephantinae, containing all other elephantids. Recent genetic research has indicated that Elephas and Mammuthus are more closely related to each other than to Loxodonta, with Palaeoloxodon closely related to Loxodonta. Palaeoloxodon also appears to have received extensive hybridisation with the African forest elephant, and to a lesser extent with mammoths.
Living species
Loxodonta (African)
L. africana African bush elephant
L. cyclotis African forest elephant
Elephas (Asiatic)
E. maximus Asian elephant
E. m. maximus Sri Lankan elephant
E. m. indicus Indian elephant
E. m. sumatranus Sumatran elephant
E. m. borneensis Borneo elephant
Classification
Elephantidae
†Stegotetrabelodon (4 species)
Subfamily Elephantinae
†Primelephas (2 species)
Elephas (7+ species)
†Stegoloxodon (2 species)
Loxodonta (6 species)
†Palaeoloxodon (14+ species)
†Phanagoroloxodon (1 species)
†Mammuthus (10 species)
†Stegodibelodon (1 species)
†Selenetherium (1 species)
Evolutionary history
Elephantids are thought to have evolved from gomphotheres, with some authors proposing the most likely ancestors to be African species of the "tetralophodont gomphothere" Tetralophodon. The earliest members of the family, are known from the Late Miocene, around 9–10 million years ago. The modern genera of elephants and mammoths had diverged from each other by the end of the Miocene, around 5 million years ago. Elephantids began to migrate out of Africa during the Pliocene, with mammoths and Elephas arriving in Eurasia around 3–3.8 million years ago. Around 1.5 million years ago, mammoths migrated into North America. At the end of the Early Pleistocene, around 0.8 million years ago, Palaeoloxodon migrated out of Africa, becoming widespread across Eurasia, from Western Europe to Japan. Palaeoloxodon and Mammuthus became extinct during the Late Pleistocene-Holocene, with the last population of mammoths persisting on Wrangel Island until around 4,000 years ago.
| Biology and health sciences | Proboscidea | Animals |
265509 | https://en.wikipedia.org/wiki/Myrtaceae | Myrtaceae | Myrtaceae (), the myrtle family, is a family of dicotyledonous plants placed within the order Myrtales. Myrtle, pōhutukawa, bay rum tree, clove, guava, acca (feijoa), allspice, and eucalyptus are some notable members of this group. All species are woody, contain essential oils, and have flower parts in multiples of four or five. The leaves are evergreen, alternate to mostly opposite, simple, and usually entire (i.e., without a toothed margin). The flowers have a base number of five petals, though in several genera, the petals are minute or absent. The stamens are usually very conspicuous, brightly coloured, and numerous.
Evolutionary history
Scientists hypothesize that the family Myrtaceae arose between 60 and 56 million years ago (Mya) during the Paleocene era. Pollen fossils have been sourced to the ancient supercontinent Gondwana. The breakup of Gondwana during the Cretaceous period (145 to 66 Mya) geographically isolated disjunct taxa and allowed for rapid speciation; in particular, genera once considered members of the now-defunct Leptospermoideae alliance are now isolated within Oceania. Generally, experts agree that vicariance is responsible for the differentiation of Myrtaceae taxa, except in the cases of Leptospermum species now located on New Zealand and New Caledonia, islands which may have been submerged at the time of late Eocene differentiation.
Diversity
Recent estimates suggest the Myrtaceae include about 5,950 species in about 132 genera. The family has a wide distribution in tropical and warm-temperate regions of the world, and is common in many of the world's biodiversity hotspots. Genera with capsular fruits such as Eucalyptus, Corymbia, Angophora, Leptospermum, and Melaleuca are absent from the Americas, apart from Metrosideros in Chile and Argentina. Genera with fleshy fruits have their greatest concentrations in eastern Australia and Malesia (the Australasian realm) and the Neotropics. Eucalyptus is a dominant, nearly ubiquitous genus in the more mesic parts of Australia and extends north sporadically to the Philippines. Eucalyptus regnans is the tallest flowering plant in the world. Other important Australian genera are Callistemon (bottlebrushes), Syzygium, and Melaleuca (paperbarks). Species of the genus Osbornia, native to Australasia, are mangroves. Eugenia, Myrcia, and Calyptranthes are among the larger genera in the neotropics.
Historically, the Myrtaceae were divided into two subfamilies. Subfamily Myrtoideae (about 75 genera) was recognized as having fleshy fruits and opposite, entire leaves. Most genera in this subfamily have one of three easily recognized types of embryos. The genera of Myrtoideae can be very difficult to distinguish in the absence of mature fruits. Myrtoideae are found worldwide in subtropical and tropical regions, with centers of diversity in the Neotropics, northeastern Australia, and Malesia. In contrast, subfamily Leptospermoideae (about 80 genera) was recognized as having dry, dehiscent fruits (capsules) and leaves arranged spirally or alternate. The Leptospermoideae are found mostly in Australasia, with a centre of diversity in Australia. Many genera in Western Australia have greatly reduced leaves and flowers typical of more xeric habitats.
Taxonomy
The division of the Myrtaceae into Leptospermoideae and Myrtoideae was challenged by a number of authors, including Johnson and Briggs (1984), who identified 14 tribes or clades within Myrtaceae, and found Myrtoideae to be polyphyletic. Molecular studies by several groups of authors, as of 2008, have confirmed the baccate (fleshy) fruits evolved twice from capsular fruits and, as such, the two-subfamily classification does not accurately portray the phylogenetic history of the family. Thus, many workers are now using a recent analysis by Wilson et al. (2001) as a starting point for further analyses of the family. This study pronounced both Leptospermoideae and Myrtoideae invalid, but retained several smaller suballiances shown to be monophyletic through matK analysis.
The genera Heteropyxis and Psiloxylon have been separated as separate families by many authors in the past as Heteropyxidaceae and Psiloxylaceae. However, Wilson et al. included them in Myrtaceae. These two genera are presently believed to be the earliest arising and surviving lineages of Myrtaceae.
The most recent classification recognizes 17 tribes and two subfamilies, Myrtoideae and Psiloxyloideae, based on a phylogenetic analysis of plastid DNA.
Many new species are being described annually from throughout the range of Myrtaceae. Likewise, new genera are being described nearly yearly.
Classification
Following Wilson (2011)
Subfamily Psiloxyloideae
tribe Psiloxyleae
tribe Heteropyxideae
Subfamily Myrtoideae
tribe Xanthostemoneae
tribe Lophostemoneae
tribe Osbornieae
tribe Melaleuceae
tribe Kanieae
tribe Backhousieae
tribe Metrosidereae
tribe Tristanieae
tribe Syzygieae
tribe Myrteae
tribe Eucalypteae
tribe Syncarpieae
tribe Lindsayomyrteae
tribe Leptospermeae
tribe Chamelaucieae
Genera
127 genera are currently accepted:
Acca
Accara
Actinodium
Agonis
Algrizea
Allosyncarpia
Aluta
Amomyrtella
Amomyrtus
Angophora
Anticoryne
Archirhodomyrtus
Arillastrum
Astartea
Asteromyrtus
Astus
Austrobaeckea
Austromyrtus
Babingtonia
Backhousia
Baeckea
Balaustion
Barongia
Basisperma
Beaufortia – synonym of Melaleuca
Blepharocalyx
Callistemon – synonym of Melaleuca
Calothamnus – synonym of Melaleuca
Calycolpus
Calycorectes
Calytrix
Campomanesia
Chamelaucium
Chamguava
Cheyniana
Cloezia
Conothamnus – synonym of Melaleuca
Corymbia
Curitiba
Cyathostemon
Darwinia
Decaspermum
Enekbatus
Eremaea – synonym of Melaleuca
Ericomyrtus
Eucalyptopsis
Eucalyptus
Eugenia
Euryomyrtus
Feijoa
Gossia
Harmogia
Heteropyxis
Homalocalyx
Homalospermum
Homoranthus
Hypocalymma
Hysterobaeckea
Kanakomyrtus
Kania
Kardomia
Kjellbergiodendron
Kunzea
Lamarchea – synonym of Melaleuca
Legrandia
Lenwebbia
Leptospermum
Lindsayomyrtus
Lithomyrtus
Lophomyrtus
Lophostemon
Luma
Lysicarpus
Malleostemon
Melaleuca
Metrosideros
Micromyrtus
Mitrantia
Mosiera
Myrceugenia
Myrcia
Myrcianthes
Myrciaria
Myrrhinium
Myrtastrum
Myrtella
Myrteola
Myrtus
Neofabricia
Neomitranthes
Neomyrtus
Nothomyrcia
Ochrosperma
Octamyrtus
Osbornia
Oxymyrrhine
†Paleomyrtinaea
Paragonis – synonym of Agonis
Pericalymma
Phymatocarpus – synonym of Melaleuca
Pileanthus
Pilidiostigma
Pimenta
Pleurocalyptus
Plinia
Psidium
Psiloxylon
Purpureostemon
Regelia – synonym of Melaleuca
Rhodamnia
Rhodomyrtus
Rinzia
Ristantia
Sannantha
Scholtzia
Seorsus
Siphoneugena
Sphaerantia
Stenostegia
Stockwellia
Syncarpia
Syzygium
Taxandria
Temu
Tetrapora
Thaleropia
Thryptomene
Triplarina
Tristania
Tristaniopsis
Ugni
Uromyrtus
Verticordia
Welchiodendron
Whiteodendron
Xanthomyrtus
Xanthostemon
Ecology
Myrtaceae is foraged by many stingless bees, especially by species such as Melipona bicolor which gather pollen from this plant family. Some Australian species such as Tetragonula hockingsi and T. carbonaria are also known to collect resin from the mature seed pods of Corymbia torelliana, resulting in mellitochory as the seeds get stuck onto the corbiculae of the bees and sometimes are successfully disposed of by colony members that remove them. But usually, they get stuck in the hives or near hive entrances instead, hence also making it a minor nuisance for some keepers as they can take up a lot of space. Fortunately, this is only known to occur in the eastern areas of Australia, but could occur in other neighbouring countries where some Corymbia species are native.
Weevils in the tribe Cryptoplini mostly use Myrtaceae as hosts. Their larvae can develop in flower and fruit buds, or in galls (often galls already formed by other insects).
| Biology and health sciences | Myrtales | null |
265533 | https://en.wikipedia.org/wiki/Juglandaceae | Juglandaceae | The Juglandaceae are a plant family known as the walnut family. They are trees, or sometimes shrubs, in the order Fagales. Members of this family are native to the Americas, Eurasia, and Southeast Asia.
The nine or ten genera in the family have a total of around 50 species, and include the commercially important nut-producing trees walnut (Juglans), pecan (Carya illinoinensis), and hickory (Carya). The Persian walnut, Juglans regia, is one of the major nut crops of the world. Walnut, hickory, and gaulin are also valuable timber trees while pecan wood is also valued as cooking fuel.
Description
Members of the walnut family have large, aromatic leaves that are usually alternate, but opposite in Alfaroa and Oreomunnea. The leaves are pinnately compound or ternate, and usually 20–100 cm long. The trees are wind-pollinated, and the flowers are usually arranged in catkins.
The fruits of the Juglandaceae are often confused with drupes but are accessory fruit because the outer covering of the fruit is technically an involucre and thus not morphologically part of the carpel; this means it cannot be a drupe but is instead a drupe-like nut.
Taxonomy
The known living genera are grouped into subfamilies, tribes, and subtribes as follows:
Subfamily Rhoipteleoideae Reveal
Rhoiptelea Diels & Hand.-Mazz. 1932
Subfamily Engelhardioideae Iljinskaya 1990
Alfaroa Standl. 1927—gaulin
Engelhardia Lesch. ex Blume 1825–1826—cheo
Oreomunnea Oerst. 1856
Subfamily Juglandoideae Eaton 1836
Tribe Platycaryeae Nakai 1933
Platycarya Siebold & Zucc. 1843
Tribe Juglandeae Rchb. 1832
Subtribe Caryinae D.E. Stone & P. S. Manos 2001
Carya Nutt. 1818—hickory and pecan
Annamocarya A.Chev. 1941 (sometimes included in Carya)
Subtribe Juglandinae D.E. Stone & P. S. Manos 2001
Cyclocarya Iljinsk 1953—wheel wingnut
Juglans L. 1753—walnut
Pterocarya Kunth 1824—wingnut
Systematics
Modern molecular phylogenetics suggest the following relationships:
| Biology and health sciences | Fagales | Plants |
265594 | https://en.wikipedia.org/wiki/Block%20and%20tackle | Block and tackle | A block and tackle or only tackle is a system of two or more pulleys with a rope or cable threaded between them, usually used to lift heavy loads.
The pulleys are assembled to form blocks and then blocks are paired so that one is fixed and one moves with the load. The rope is threaded through the pulleys to provide mechanical advantage that amplifies the force applied to the rope.
Hero of Alexandria described cranes formed from assemblies of pulleys in the first century. Illustrated versions of Hero's Mechanica (a book on raising heavy weights) show early block and tackle systems.
Overview
A block is a set of pulleys or sheaves mounted on a single frame. An assembly of blocks with a rope threaded through the pulleys is called tackle. The process of threading ropes or cables through blocks is called "reeving", and a threaded block and tackle is said to have been "rove". A block and tackle system amplifies the tension force in the rope to lift heavy loads. They are common on boats and sailing ships, where tasks are often performed manually, as well as on cranes and drilling rigs, where once rove, the tasks are performed by heavy equipment.
In the diagram shown here, the number of rope sections of the tackles shown is as follows:
Gun tackle: 2
Luff tackle: 3
Double tackle: 4
Gyn tackle: 5
Threefold purchase: 6
Note that in the gun tackle, double tackle and threefold purchase, both blocks have the same number of pulleys (one, two and three, respectively) whereas the Luff tackle and Gyn tackle have mis-matched blocks with differing numbers of pulleys.
Mechanical advantage
A block and tackle is characterized by the use of a single continuous rope to transmit a tension force around one or more pulleys to lift or move a load. Its mechanical advantage is the number of parts of the rope that act on the load. The mechanical advantage of a tackle dictates how much easier it is to haul or lift the load.
If frictional losses are neglected, the mechanical advantage of a block and tackle is equal to the number of parts in the line that either attach to or run through the moving blocks—in other words, the number of supporting rope sections.
An ideal block and tackle with a moving block supported by n rope sections has the mechanical advantage (MA),
where FA is the hauling (or input) force and FB is the load.
Consider the set of pulleys that form the moving block and the parts of the rope that support this block. If there are n of these parts of the rope supporting the load FB, then a force balance on the moving block shows that the tension in each of the parts of the rope must be FB/n. This means the input force on the rope is FA=FB/n. Thus, the block and tackle reduces the input force by the factor n.
Ideal mechanical advantage correlates directly with velocity ratio. The velocity ratio of a tackle is the ratio between the velocity of the hauling line to that of the hauled load. A line with a mechanical advantage of 4 has a velocity ratio of 4:1. In other words, to raise a load at 1 metre per second, the hauling part of the rope must be pulled at 4 metres per second. Therefore, the mechanical advantage of a double tackle is 4.
Rove to (dis)advantage
The mechanical advantage of any tackle can be increased by interchanging the fixed and moving blocks so the rope is attached to the moving block and the rope is pulled in the direction of the lifted load. In this case the block and tackle is said to be "rove to advantage."
"Rove to advantage" – where the pull on the rope is in the same direction as that in which the load is to be moved. The hauling part is pulled from the moving block.
"Rove to disadvantage" – where the pull on the rope is in the opposite direction to that in which the load is to be moved. The hauling part is pulled from the fixed block.
Diagram 3 shows three rope parts supporting the load W, which means the tension in the rope is W/3. Thus, the mechanical advantage is three-to-one.
By adding a pulley to the fixed block of a gun tackle the direction of the pulling force is reversed though the mechanical advantage remains the same, Diagram 3a. This is an example of the Luff tackle.
The decision of which to use depends on pragmatic considerations for the total ergonomics of working with a particular situation. Reeving to advantage is the most efficient use of equipment and resources. For example, if the load is to be hauled parallel to the ground, reeving to advantage enables the pulling force to be in the direction of the load movement, allowing obstacles to be managed more easily.
Reeving to disadvantage adds an extra sheave to change the direction of the pulling line to a potentially more ergonomic direction, which increases friction losses without improving the velocity ratio. Situations in which reeving to disadvantage may be more desirable include lifting from a fixed point overhead--the additional pulley allows pulling downwards instead of upwards so that the weight of the lifter can offset the weight of the load, or allows pulling sideways, enabling multiple lifters to combine effort.
Friction
The formula used to find the effort required to raise a given weight using a block and fall:
where is the force applied to the hauling part of the line (the input force), is the weight of the load (the output force), is the ideal mechanical advantage of the system (which is the same as the number of segments of line extending from the moving block), and is the mechanical efficiency of the system (equal to one for an ideal frictionless system; a fraction less than one for real-world systems with energy losses due to friction and other causes). If is the number of sheaves in the purchase, and there is a roughly % loss of efficiency at each sheave due to friction, then:
This approximation is more accurate for smaller values of and . A more precise estimate of efficiency is possible by use of the sheave friction factor, (which may be obtainable from the manufacturer or published tables). The relevant equation is:
Typical values are 1.04 for roller bearing sheaves and 1.09 for plain bearing sheaves (with wire rope).
The increased force produced by a tackle is offset by both the increased length of rope needed and the friction in the system. In order to raise a block and tackle with a mechanical advantage of 6 a distance of 1 metre, it is necessary to pull 6 metres of rope through the blocks. Frictional losses also mean there is a practical point at which the benefit of adding a further sheave is offset by the incremental increase in friction which would require additional force to be applied in order to lift the load. Too much friction may result in the tackle not allowing the load to be released easily, or by the reduction in force needed to move the load being judged insufficient because undue friction has to be overcome as well.
Mid-line attachment
When installing a block on an existing line, it is often inconvenient at best to thread the rope through the block to be added.
Open blocks have a space wide enough between the fixed cheeks to be able to slide the pulley over the rope. These can be extremely small and light while retaining significant strength due to the lack of moving parts.
A swing cheek block is a special kind of block which can be opened to engage with a bight, without the necessity to thread the rope through the block or remove the load from the end of the rope. The snatch block is also the load lifting pulley in certain arrangements, such as during use in a recovery operation.
Swing cheek blocks may be roughly divided into two categories:
Swing cheek pulleys: used for light loads or redirection of forces, usually with a single pulley wheel (though multiple sheaves/cheeks are not uncommon) and an attachment point (or several) for a carabiner or sling. The cheeks are not fixed or locked in position aside from the device used to secure them to the load or rigging point.
Examples of use (in an arboricultural setting) include: tail minding/tending, and for setting a rigging point in the tree above the cut to take place—a positive rigging situation.
Snatch or impact blocks: used for heavier loads and more dynamic rigging, the cheeks of these blocks are fixed in place with a pin which locks into the opposite cheek. This pin may form part of the axle for a second pulley, which is secured to the load or rigging point with a soft sling, rather than a solid device such as a shackle. This allows for more even distribution of forces to the faces where the forces will be applied, as opposed to a carabiner or shackle, where the forces are applied more strongly to corners and edges, increasing risk of deformation or damage.
Examples of use (again, in relation to tree care) may include setting a block below the current cut, resulting in a 'negative' rigging situation, in which shock loads can be significant—especially if removing large sections of vertical stem.
Literature
M. Oppolzer, T. Wahls: I Like To Move It. Flaschenzüge in der Seiltechnik. Hamburg 2019, .
Rescue Technician: Operational Readiness for Rescue Providers. St. Louis, Missouri 1998, .
| Technology | Rigid components | null |
265787 | https://en.wikipedia.org/wiki/Tyrant%20flycatcher | Tyrant flycatcher | The tyrant flycatchers (Tyrannidae) are a family of passerine birds which occur throughout North and South America. They are the largest family of birds, with more than 400 species. They are the most diverse avian family in every country in the Americas, except for the United States and Canada. The members vary greatly in shape, patterns, size and colors. Some tyrant flycatchers may superficially resemble the Old World flycatchers, which they are named after but are not closely related to. They are members of suborder Tyranni (suboscines), which do not have the sophisticated vocal capabilities of most other songbirds.
A number of species previously included in this family are now placed in the family Tityridae (see Systematics). Sibley and Alquist in their 1990 bird taxonomy had the genera Mionectes, Leptopogon, Pseudotriccus, Poecilotriccus, Taenotriccus, Hemitriccus, Todirostrum and Corythopis as a separate family Pipromorphidae, but although it is still thought that these genera are basal to most of the family, they are not each other's closest relatives.
Description
Most species are rather plain, with various hues of brown, gray and white commonplace, often providing some degree of presumed camouflage. Obvious exceptions include the bright red vermilion flycatcher, blue, black, white and yellow many-colored rush-tyrant and some species of tody-flycatchers or tyrants, which are often yellow, black, white and/or rufous, from the Todirostrum, Hemitriccus and Poecilotriccus genera. Several species have bright yellow underparts, from the ornate flycatcher to the great kiskadee. Some species have erectile crests. Several of the large genera (i.e. Elaenia, Myiarchus or Empidonax) are quite difficult to tell apart in the field due to similar plumage and some are best distinguished by their voices. Behaviorally they can vary from species such as spadebills which are tiny, shy and live in dense forest interiors to kingbirds, which are relatively large, bold, inquisitive and often inhabit open areas near human habitations. As the name implies, a great majority of tyrant flycatchers are almost entirely insectivorous (though not necessarily specialized in flies). Tyrant flycatchers are largely opportunistic feeders and often catch any flying or arboreal insect they encounter. However, food can vary greatly and some (like the large great kiskadee) will eat fruit or small vertebrates (e.g. small frogs). In North America, most species are associated with a "sallying" feeding style, where they fly up to catch an insect directly from their perch and then immediately return to the same perch. Most tropical species, however, do not feed in this fashion and several types prefer to glean insects from leaves and bark. Tropical species are sometimes found in mixed-species foraging flocks, where various types of passerines and other smallish birds are found feeding in proximity.
The smallest family members are the closely related short-tailed pygmy tyrant and black-capped pygmy tyrant from the genus Myiornis (the first species usually being considered marginally smaller on average). These species reach a total length of and a weight of . By length, they are the smallest passerines on earth, although some species of Old World warblers apparently rival them in their minuscule mean body masses if not in total length. The minuscule size and very short tail of the Myiornis pygmy tyrants often lend them a resemblance to a tiny ball or insect. The largest tyrant flycatcher is the great shrike-tyrant at and . A few species such as the streamer-tailed tyrant, scissor-tailed flycatcher and fork-tailed flycatcher have a larger total length — up to in the fork-tailed flycatcher at least — but this is mainly due to their extremely long tails; the fork-tailed flycatcher has the longest tail feathers of any known bird relative to their size (this being in reference to true tail feathers, not to be confused with elongated tail streamers as seen in some from the Phasianidae family of galliforms).
Habitat and distribution
Species richness of Tyrannidae, when compared to habitat, is highly variable, although most every land habitat in the Americas has at least some of these birds. The habitats of tropical lowland evergreen forest and montane evergreen forest have the highest single site species diversity while many habitats including rivers, palm forest, white sand forest, tropical deciduous forest edge, southern temperate forest, southern temperate forest edge, semi-humid/humid montane scrub, and northern temperate grassland have the lowest single species diversity. The variation between the highest and the lowest is extreme; ninety species can be found in the tropical lowland evergreen forests while the number of species that can be found in the habitats listed above typically are in the single digits. This may be due in part to the fewer niches found in certain areas and therefore fewer places for the species to occupy.
Tyrannidae specialization among habitats is very strong in tropical lowland evergreen forests and montane evergreen forests. These habitat types, therefore, display the greatest specialization. The counts differ by three species (tropical lowland evergreen forests have 49 endemic species and montane evergreen forests have 46 endemic species). It can be assumed that they both have similar levels of specialization.
Regionally, the Atlantic Forest has the highest species richness with the Chocó following closely behind.
Status and conservation
The northern beardless tyrannulet (Camptostoma imberbe) is protected under the Migratory Bird Treaty Act of 1918. This species is common south of the US border. The situation for a number of other species from South and Central America is far more problematic. In 2007, BirdLife International (and consequently IUCN) considered two species, the Minas Gerais tyrannulet and Kaempfer's tody-tyrant critically endangered. Both are endemic to Brazil. Additionally, seven species were considered endangered and eighteen species vulnerable.
Systematics
The family's name is derived from an early description of the eastern kingbird as "the tyrant" by naturalist Mark Catesby in the 1730s. Carl Linnaeus adopted that name for the entire family Tyrannidae, because he admired Catesby's work.
The family contains 447 species divided into 104 genera. A full list, sortable by common and binomial names, is at list of tyrant flycatcher species. Species in the genera Tityra, Pachyramphus, Laniocera and Xenopsaris were formerly placed in this family, but evidence suggested they belong in their own family, the Tityridae, where they are now placed by SACC.
| Biology and health sciences | Tyranni | null |
265823 | https://en.wikipedia.org/wiki/Thermodynamic%20equilibrium | Thermodynamic equilibrium | Thermodynamic equilibrium is a notion of thermodynamics with axiomatic status referring to an internal state of a single thermodynamic system, or a relation between several thermodynamic systems connected by more or less permeable or impermeable walls. In thermodynamic equilibrium, there are no net macroscopic flows of mass nor of energy within a system or between systems. In a system that is in its own state of internal thermodynamic equilibrium, not only is there an absence of macroscopic change, but there is an “absence of any tendency toward change on a macroscopic scale.”
Systems in mutual thermodynamic equilibrium are simultaneously in mutual thermal, mechanical, chemical, and radiative equilibria. Systems can be in one kind of mutual equilibrium, while not in others. In thermodynamic equilibrium, all kinds of equilibrium hold at once and indefinitely, unless disturbed by a thermodynamic operation. In a macroscopic equilibrium, perfectly or almost perfectly balanced microscopic exchanges occur; this is the physical explanation of the notion of macroscopic equilibrium.
A thermodynamic system in a state of internal thermodynamic equilibrium has a spatially uniform temperature. Its intensive properties, other than temperature, may be driven to spatial inhomogeneity by an unchanging long-range force field imposed on it by its surroundings.
In systems that are at a state of non-equilibrium there are, by contrast, net flows of matter or energy. If such changes can be triggered to occur in a system in which they are not already occurring, the system is said to be in a "meta-stable equilibrium".
Though not a widely named "law," it is an axiom of thermodynamics that there exist states of thermodynamic equilibrium. The second law of thermodynamics states that when an isolated body of material starts from an equilibrium state, in which portions of it are held at different states by more or less permeable or impermeable partitions, and a thermodynamic operation removes or makes the partitions more permeable, then it spontaneously reaches its own new state of internal thermodynamic equilibrium and this is accompanied by an increase in the sum of the entropies of the portions.
Overview
Classical thermodynamics deals with states of dynamic equilibrium. The state of a system at thermodynamic equilibrium is the one for which some thermodynamic potential is minimized (in the absence of an applied voltage), or for which the entropy (S) is maximized, for specified conditions. One such potential is the Helmholtz free energy (A), for a closed system at constant volume and temperature (controlled by a heat bath):
Another potential, the Gibbs free energy (G), is minimized at thermodynamic equilibrium in a closed system at constant temperature and pressure, both controlled by the surroundings:
where T denotes the absolute thermodynamic temperature, P the pressure, S the entropy, V the volume, and U the internal energy of the system. In other words, is a necessary condition for chemical equilibrium under these conditions (in the absence of an applied voltage).
Thermodynamic equilibrium is the unique stable stationary state that is approached or eventually reached as the system interacts with its surroundings over a long time. The above-mentioned potentials are mathematically constructed to be the thermodynamic quantities that are minimized under the particular conditions in the specified surroundings.
Conditions
For a completely isolated system, S is maximum at thermodynamic equilibrium.
For a closed system at controlled constant temperature and volume, A is minimum at thermodynamic equilibrium.
For a closed system at controlled constant temperature and pressure without an applied voltage, G is minimum at thermodynamic equilibrium.
The various types of equilibriums are achieved as follows:
Two systems are in thermal equilibrium when their temperatures are the same.
Two systems are in mechanical equilibrium when their pressures are the same.
Two systems are in diffusive equilibrium when their chemical potentials are the same.
All forces are balanced and there is no significant external driving force.
Relation of exchange equilibrium between systems
Often the surroundings of a thermodynamic system may also be regarded as another thermodynamic system. In this view, one may consider the system and its surroundings as two systems in mutual contact, with long-range forces also linking them. The enclosure of the system is the surface of contiguity or boundary between the two systems. In the thermodynamic formalism, that surface is regarded as having specific properties of permeability. For example, the surface of contiguity may be supposed to be permeable only to heat, allowing energy to transfer only as heat. Then the two systems are said to be in thermal equilibrium when the long-range forces are unchanging in time and the transfer of energy as heat between them has slowed and eventually stopped permanently; this is an example of a contact equilibrium. Other kinds of contact equilibrium are defined by other kinds of specific permeability. When two systems are in contact equilibrium with respect to a particular kind of permeability, they have common values of the intensive variable that belongs to that particular kind of permeability. Examples of such intensive variables are temperature, pressure, chemical potential.
A contact equilibrium may be regarded also as an exchange equilibrium. There is a zero balance of rate of transfer of some quantity between the two systems in contact equilibrium. For example, for a wall permeable only to heat, the rates of diffusion of internal energy as heat between the two systems are equal and opposite. An adiabatic wall between the two systems is 'permeable' only to energy transferred as work; at mechanical equilibrium the rates of transfer of energy as work between them are equal and opposite. If the wall is a simple wall, then the rates of transfer of volume across it are also equal and opposite; and the pressures on either side of it are equal. If the adiabatic wall is more complicated, with a sort of leverage, having an area-ratio, then the pressures of the two systems in exchange equilibrium are in the inverse ratio of the volume exchange ratio; this keeps the zero balance of rates of transfer as work.
A radiative exchange can occur between two otherwise separate systems. Radiative exchange equilibrium prevails when the two systems have the same temperature.
Thermodynamic state of internal equilibrium of a system
A collection of matter may be entirely isolated from its surroundings. If it has been left undisturbed for an indefinitely long time, classical thermodynamics postulates that it is in a state in which no changes occur within it, and there are no flows within it. This is a thermodynamic state of internal equilibrium. (This postulate is sometimes, but not often, called the "minus first" law of thermodynamics. One textbook calls it the "zeroth law", remarking that the authors think this more befitting that title than its more customary definition, which apparently was suggested by Fowler.)
Such states are a principal concern in what is known as classical or equilibrium thermodynamics, for they are the only states of the system that are regarded as well defined in that subject. A system in contact equilibrium with another system can by a thermodynamic operation be isolated, and upon the event of isolation, no change occurs in it. A system in a relation of contact equilibrium with another system may thus also be regarded as being in its own state of internal thermodynamic equilibrium.
Multiple contact equilibrium
The thermodynamic formalism allows that a system may have contact with several other systems at once, which may or may not also have mutual contact, the contacts having respectively different permeabilities. If these systems are all jointly isolated from the rest of the world those of them that are in contact then reach respective contact equilibria with one another.
If several systems are free of adiabatic walls between each other, but are jointly isolated from the rest of the world, then they reach a state of multiple contact equilibrium, and they have a common temperature, a total internal energy, and a total entropy. Amongst intensive variables, this is a unique property of temperature. It holds even in the presence of long-range forces. (That is, there is no "force" that can maintain temperature discrepancies.) For example, in a system in thermodynamic equilibrium in a vertical gravitational field, the pressure on the top wall is less than that on the bottom wall, but the temperature is the same everywhere.
A thermodynamic operation may occur as an event restricted to the walls that are within the surroundings, directly affecting neither the walls of contact of the system of interest with its surroundings, nor its interior, and occurring within a definitely limited time. For example, an immovable adiabatic wall may be placed or removed within the surroundings. Consequent upon such an operation restricted to the surroundings, the system may be for a time driven away from its own initial internal state of thermodynamic equilibrium. Then, according to the second law of thermodynamics, the whole undergoes changes and eventually reaches a new and final equilibrium with the surroundings. Following Planck, this consequent train of events is called a natural thermodynamic process. It is allowed in equilibrium thermodynamics just because the initial and final states are of thermodynamic equilibrium, even though during the process there is transient departure from thermodynamic equilibrium, when neither the system nor its surroundings are in well defined states of internal equilibrium. A natural process proceeds at a finite rate for the main part of its course. It is thereby radically different from a fictive quasi-static 'process' that proceeds infinitely slowly throughout its course, and is fictively 'reversible'. Classical thermodynamics allows that even though a process may take a very long time to settle to thermodynamic equilibrium, if the main part of its course is at a finite rate, then it is considered to be natural, and to be subject to the second law of thermodynamics, and thereby irreversible. Engineered machines and artificial devices and manipulations are permitted within the surroundings. The allowance of such operations and devices in the surroundings but not in the system is the reason why Kelvin in one of his statements of the second law of thermodynamics spoke of "inanimate" agency; a system in thermodynamic equilibrium is inanimate.
Otherwise, a thermodynamic operation may directly affect a wall of the system.
It is often convenient to suppose that some of the surrounding subsystems are so much larger than the system that the process can affect the intensive variables only of the surrounding subsystems, and they are then called reservoirs for relevant intensive variables.
Local and global equilibrium
It can be useful to distinguish between global and local thermodynamic equilibrium. In thermodynamics, exchanges within a system and between the system and the outside are controlled by intensive parameters. As an example, temperature controls heat exchanges. Global thermodynamic equilibrium (GTE) means that those intensive parameters are homogeneous throughout the whole system, while local thermodynamic equilibrium (LTE) means that those intensive parameters are varying in space and time, but are varying so slowly that, for any point, one can assume thermodynamic equilibrium in some neighborhood about that point.
If the description of the system requires variations in the intensive parameters that are too large, the very assumptions upon which the definitions of these intensive parameters are based will break down, and the system will be in neither global nor local equilibrium. For example, it takes a certain number of collisions for a particle to equilibrate to its surroundings. If the average distance it has moved during these collisions removes it from the neighborhood it is equilibrating to, it will never equilibrate, and there will be no LTE. Temperature is, by definition, proportional to the average internal energy of an equilibrated neighborhood. Since there is no equilibrated neighborhood, the concept of temperature doesn't hold, and the temperature becomes undefined.
This local equilibrium may apply only to a certain subset of particles in the system. For example, LTE is usually applied only to massive particles. In a radiating gas, the photons being emitted and absorbed by the gas do not need to be in a thermodynamic equilibrium with each other or with the massive particles of the gas for LTE to exist. In some cases, it is not considered necessary for free electrons to be in equilibrium with the much more massive atoms or molecules for LTE to exist.
As an example, LTE will exist in a glass of water that contains a melting ice cube. The temperature inside the glass can be defined at any point, but it is colder near the ice cube than far away from it. If energies of the molecules located near a given point are observed, they will be distributed according to the Maxwell–Boltzmann distribution for a certain temperature. If the energies of the molecules located near another point are observed, they will be distributed according to the Maxwell–Boltzmann distribution for another temperature.
Local thermodynamic equilibrium does not require either local or global stationarity. In other words, each small locality need not have a constant temperature. However, it does require that each small locality change slowly enough to practically sustain its local Maxwell–Boltzmann distribution of molecular velocities. A global non-equilibrium state can be stably stationary only if it is maintained by exchanges between the system and the outside. For example, a globally-stable stationary state could be maintained inside the glass of water by continuously adding finely powdered ice into it to compensate for the melting, and continuously draining off the meltwater. Natural transport phenomena may lead a system from local to global thermodynamic equilibrium. Going back to our example, the diffusion of heat will lead our glass of water toward global thermodynamic equilibrium, a state in which the temperature of the glass is completely homogeneous.
Reservations
Careful and well informed writers about thermodynamics, in their accounts of thermodynamic equilibrium, often enough make provisos or reservations to their statements. Some writers leave such reservations merely implied or more or less unstated.
For example, one widely cited writer, H. B. Callen writes in this context: "In actuality, few systems are in absolute and true equilibrium." He refers to radioactive processes and remarks that they may take "cosmic times to complete, [and] generally can be ignored". He adds "In practice, the criterion for equilibrium is circular. Operationally, a system is in an equilibrium state if its properties are consistently described by thermodynamic theory!"
J.A. Beattie and I. Oppenheim write: "Insistence on a strict interpretation of the definition of equilibrium would rule out the application of thermodynamics to practically all states of real systems."
Another author, cited by Callen as giving a "scholarly and rigorous treatment", and cited by Adkins as having written a "classic text", A.B. Pippard writes in that text: "Given long enough a supercooled vapour will eventually condense, ... . The time involved may be so enormous, however, perhaps 10100 years or more, ... . For most purposes, provided the rapid change is not artificially stimulated, the systems may be regarded as being in equilibrium."
Another author, A. Münster, writes in this context. He observes that thermonuclear processes often occur so slowly that they can be ignored in thermodynamics. He comments: "The concept 'absolute equilibrium' or 'equilibrium with respect to all imaginable processes', has therefore, no physical significance." He therefore states that: "... we can consider an equilibrium only with respect to specified processes and defined experimental conditions."
According to L. Tisza: "... in the discussion of phenomena near absolute zero. The absolute predictions of the classical theory become particularly vague because the occurrence of frozen-in nonequilibrium states is very common."
Definitions
The most general kind of thermodynamic equilibrium of a system is through contact with the surroundings that allows simultaneous passages of all chemical substances and all kinds of energy. A system in thermodynamic equilibrium may move with uniform acceleration through space but must not change its shape or size while doing so; thus it is defined by a rigid volume in space. It may lie within external fields of force, determined by external factors of far greater extent than the system itself, so that events within the system cannot in an appreciable amount affect the external fields of force. The system can be in thermodynamic equilibrium only if the external force fields are uniform, and are determining its uniform acceleration, or if it lies in a non-uniform force field but is held stationary there by local forces, such as mechanical pressures, on its surface.
Thermodynamic equilibrium is a primitive notion of the theory of thermodynamics. According to P.M. Morse: "It should be emphasized that the fact that there are thermodynamic states, ..., and the fact that there are thermodynamic variables which are uniquely specified by the equilibrium state ... are not conclusions deduced logically from some philosophical first principles. They are conclusions ineluctably drawn from more than two centuries of experiments." This means that thermodynamic equilibrium is not to be defined solely in terms of other theoretical concepts of thermodynamics. M. Bailyn proposes a fundamental law of thermodynamics that defines and postulates the existence of states of thermodynamic equilibrium.
Textbook definitions of thermodynamic equilibrium are often stated carefully, with some reservation or other.
For example, A. Münster writes: "An isolated system is in thermodynamic equilibrium when, in the system, no changes of state are occurring at a measurable rate." There are two reservations stated here; the system is isolated; any changes of state are immeasurably slow. He discusses the second proviso by giving an account of a mixture oxygen and hydrogen at room temperature in the absence of a catalyst. Münster points out that a thermodynamic equilibrium state is described by fewer macroscopic variables than is any other state of a given system. This is partly, but not entirely, because all flows within and through the system are zero.
R. Haase's presentation of thermodynamics does not start with a restriction to thermodynamic equilibrium because he intends to allow for non-equilibrium thermodynamics. He considers an arbitrary system with time invariant properties. He tests it for thermodynamic equilibrium by cutting it off from all external influences, except external force fields. If after insulation, nothing changes, he says that the system was in equilibrium.
In a section headed "Thermodynamic equilibrium", H.B. Callen defines equilibrium states in a paragraph. He points out that they "are determined by intrinsic factors" within the system. They are "terminal states", towards which the systems evolve, over time, which may occur with "glacial slowness". This statement does not explicitly say that for thermodynamic equilibrium, the system must be isolated; Callen does not spell out what he means by the words "intrinsic factors".
Another textbook writer, C.J. Adkins, explicitly allows thermodynamic equilibrium to occur in a system which is not isolated. His system is, however, closed with respect to transfer of matter. He writes: "In general, the approach to thermodynamic equilibrium will involve both thermal and work-like interactions with the surroundings." He distinguishes such thermodynamic equilibrium from thermal equilibrium, in which only thermal contact is mediating transfer of energy.
Another textbook author, J.R. Partington, writes: "(i) An equilibrium state is one which is independent of time." But, referring to systems "which are only apparently in equilibrium", he adds : "Such systems are in states of ″false equilibrium.″" Partington's statement does not explicitly state that the equilibrium refers to an isolated system. Like Münster, Partington also refers to the mixture of oxygen and hydrogen. He adds a proviso that "In a true equilibrium state, the smallest change of any external condition which influences the state will produce a small change of state ..." This proviso means that thermodynamic equilibrium must be stable against small perturbations; this requirement is essential for the strict meaning of thermodynamic equilibrium.
A student textbook by F.H. Crawford has a section headed "Thermodynamic Equilibrium". It distinguishes several drivers of flows, and then says: "These are examples of the apparently universal tendency of isolated systems toward a state of complete mechanical, thermal, chemical, and electrical—or, in a single word, thermodynamic—equilibrium."
A monograph on classical thermodynamics by H.A. Buchdahl considers the "equilibrium of a thermodynamic system", without actually writing the phrase "thermodynamic equilibrium". Referring to systems closed to exchange of matter, Buchdahl writes: "If a system is in a terminal condition which is properly static, it will be said to be in equilibrium." Buchdahl's monograph also discusses amorphous glass, for the purposes of thermodynamic description. It states: "More precisely, the glass may be regarded as being in equilibrium so long as experimental tests show that 'slow' transitions are in effect reversible." It is not customary to make this proviso part of the definition of thermodynamic equilibrium, but the converse is usually assumed: that if a body in thermodynamic equilibrium is subject to a sufficiently slow process, that process may be considered to be sufficiently nearly reversible, and the body remains sufficiently nearly in thermodynamic equilibrium during the process.
A. Münster carefully extends his definition of thermodynamic equilibrium for isolated systems by introducing a concept of contact equilibrium. This specifies particular processes that are allowed when considering thermodynamic equilibrium for non-isolated systems, with special concern for open systems, which may gain or lose matter from or to their surroundings. A contact equilibrium is between the system of interest and a system in the surroundings, brought into contact with the system of interest, the contact being through a special kind of wall; for the rest, the whole joint system is isolated. Walls of this special kind were also considered by C. Carathéodory, and are mentioned by other writers also. They are selectively permeable. They may be permeable only to mechanical work, or only to heat, or only to some particular chemical substance. Each contact equilibrium defines an intensive parameter; for example, a wall permeable only to heat defines an empirical temperature. A contact equilibrium can exist for each chemical constituent of the system of interest. In a contact equilibrium, despite the possible exchange through the selectively permeable wall, the system of interest is changeless, as if it were in isolated thermodynamic equilibrium. This scheme follows the general rule that "... we can consider an equilibrium only with respect to specified processes and defined experimental conditions." Thermodynamic equilibrium for an open system means that, with respect to every relevant kind of selectively permeable wall, contact equilibrium exists when the respective intensive parameters of the system and surroundings are equal. This definition does not consider the most general kind of thermodynamic equilibrium, which is through unselective contacts. This definition does not simply state that no current of matter or energy exists in the interior or at the boundaries; but it is compatible with the following definition, which does so state.
M. Zemansky also distinguishes mechanical, chemical, and thermal equilibrium. He then writes: "When the conditions for all three types of equilibrium are satisfied, the system is said to be in a state of thermodynamic equilibrium".
P.M. Morse writes that thermodynamics is concerned with "states of thermodynamic equilibrium". He also uses the phrase "thermal equilibrium" while discussing transfer of energy as heat between a body and a heat reservoir in its surroundings, though not explicitly defining a special term 'thermal equilibrium'.
J.R. Waldram writes of "a definite thermodynamic state". He defines the term "thermal equilibrium" for a system "when its observables have ceased to change over time". But shortly below that definition he writes of a piece of glass that has not yet reached its "full thermodynamic equilibrium state".
Considering equilibrium states, M. Bailyn writes: "Each intensive variable has its own type of equilibrium." He then defines thermal equilibrium, mechanical equilibrium, and material equilibrium. Accordingly, he writes: "If all the intensive variables become uniform, thermodynamic equilibrium is said to exist." He is not here considering the presence of an external force field.
J.G. Kirkwood and I. Oppenheim define thermodynamic equilibrium as follows: "A system is in a state of thermodynamic equilibrium if, during the time period allotted for experimentation, (a) its intensive properties are independent of time and (b) no current of matter or energy exists in its interior or at its boundaries with the surroundings." It is evident that they are not restricting the definition to isolated or to closed systems. They do not discuss the possibility of changes that occur with "glacial slowness", and proceed beyond the time period allotted for experimentation. They note that for two systems in contact, there exists a small subclass of intensive properties such that if all those of that small subclass are respectively equal, then all respective intensive properties are equal. States of thermodynamic equilibrium may be defined by this subclass, provided some other conditions are satisfied.
Characteristics of a state of internal thermodynamic equilibrium
Homogeneity in the absence of external forces
A thermodynamic system consisting of a single phase in the absence of external forces, in its own internal thermodynamic equilibrium, is homogeneous. This means that the material in any small volume element of the system can be interchanged with the material of any other geometrically congruent volume element of the system, and the effect is to leave the system thermodynamically unchanged. In general, a strong external force field makes a system of a single phase in its own internal thermodynamic equilibrium inhomogeneous with respect to some intensive variables. For example, a relatively dense component of a mixture can be concentrated by centrifugation.
Uniform temperature
Such equilibrium inhomogeneity, induced by external forces, does not occur for the intensive variable temperature. According to E.A. Guggenheim, "The most important conception of thermodynamics is temperature." Planck introduces his treatise with a brief account of heat and temperature and thermal equilibrium, and then announces: "In the following we shall deal chiefly with homogeneous, isotropic bodies of any form, possessing throughout their substance the same temperature and density, and subject to a uniform pressure acting everywhere perpendicular to the surface." As did Carathéodory, Planck was setting aside surface effects and external fields and anisotropic crystals. Though referring to temperature, Planck did not there explicitly refer to the concept of thermodynamic equilibrium. In contrast, Carathéodory's scheme of presentation of classical thermodynamics for closed systems postulates the concept of an "equilibrium state" following Gibbs (Gibbs speaks routinely of a "thermodynamic state"), though not explicitly using the phrase 'thermodynamic equilibrium', nor explicitly postulating the existence of a temperature to define it.
Although thermodynamic laws are immutable, systems can be created that delay the time to reach thermodynamic equilibrium. In a thought experiment, Reed A. Howald conceived of a system called "The Fizz Keeper"consisting of a cap with a nozzle that can re-pressurize any standard bottle of carbonated beverage. Nitrogen and oxygen, which air are mostly made out of, would keep getting pumped in, which would slow down the rate at which the carbon dioxide fizzles out of the system. This is possible because the thermodynamic equilibrium between the unconverted and converted carbon dioxide inside the bottle would stay the same. To come to this conclusion, he also appeals to Henry's Law, which states that gases dissolve in direct proportion to their partial pressures. By influencing the partial pressure on the top of a closed system, this would help slow down the rate of fizzing out of carbonated beverages which is governed by thermodynamic equilibrium. The equilibria of carbon dioxide and other gases would not change, however the partial pressure on top would slow down the rate of dissolution extending the time a gas stays in a particular state. due to the nature of thermal equilibrium of the remainder of the beverage. The equilibrium constant of carbon dioxide would be completely independent of the nitrogen and oxygen pumped into the system, which would slow down the diffusion of gas, and yet not have an impact on the thermodynamics of the entire system.
The temperature within a system in thermodynamic equilibrium is uniform in space as well as in time. In a system in its own state of internal thermodynamic equilibrium, there are no net internal macroscopic flows. In particular, this means that all local parts of the system are in mutual radiative exchange equilibrium. This means that the temperature of the system is spatially uniform. This is so in all cases, including those of non-uniform external force fields. For an externally imposed gravitational field, this may be proved in macroscopic thermodynamic terms, by the calculus of variations, using the method of Langrangian multipliers. Considerations of kinetic theory or statistical mechanics also support this statement.
In order that a system may be in its own internal state of thermodynamic equilibrium, it is of course necessary, but not sufficient, that it be in its own internal state of thermal equilibrium; it is possible for a system to reach internal mechanical equilibrium before it reaches internal thermal equilibrium.
Number of real variables needed for specification
In his exposition of his scheme of closed system equilibrium thermodynamics, C. Carathéodory initially postulates that experiment reveals that a definite number of real variables define the states that are the points of the manifold of equilibria. In the words of Prigogine and Defay (1945): "It is a matter of experience that when we have specified a certain number of macroscopic properties of a system, then all the other properties are fixed." As noted above, according to A. Münster, the number of variables needed to define a thermodynamic equilibrium is the least for any state of a given isolated system. As noted above, J.G. Kirkwood and I. Oppenheim point out that a state of thermodynamic equilibrium may be defined by a special subclass of intensive variables, with a definite number of members in that subclass.
If the thermodynamic equilibrium lies in an external force field, it is only the temperature that can in general be expected to be spatially uniform. Intensive variables other than temperature will in general be non-uniform if the external force field is non-zero. In such a case, in general, additional variables are needed to describe the spatial non-uniformity.
Stability against small perturbations
As noted above, J.R. Partington points out that a state of thermodynamic equilibrium is stable against small transient perturbations. Without this condition, in general, experiments intended to study systems in thermodynamic equilibrium are in severe difficulties.
Approach to thermodynamic equilibrium within an isolated system
When a body of material starts from a non-equilibrium state of inhomogeneity or chemical non-equilibrium, and is then isolated, it spontaneously evolves towards its own internal state of thermodynamic equilibrium. It is not necessary that all aspects of internal thermodynamic equilibrium be reached simultaneously; some can be established before others. For example, in many cases of such evolution, internal mechanical equilibrium is established much more rapidly than the other aspects of the eventual thermodynamic equilibrium. Another example is that, in many cases of such evolution, thermal equilibrium is reached much more rapidly than chemical equilibrium.
Fluctuations within an isolated system in its own internal thermodynamic equilibrium
In an isolated system, thermodynamic equilibrium by definition persists over an indefinitely long time. In classical physics it is often convenient to ignore the effects of measurement and this is assumed in the present account.
To consider the notion of fluctuations in an isolated thermodynamic system, a convenient example is a system specified by its extensive state variables, internal energy, volume, and mass composition. By definition they are time-invariant. By definition, they combine with time-invariant nominal values of their conjugate intensive functions of state, inverse temperature, pressure divided by temperature, and the chemical potentials divided by temperature, so as to exactly obey the laws of thermodynamics. But the laws of thermodynamics, combined with the values of the specifying extensive variables of state, are not sufficient to provide knowledge of those nominal values. Further information is needed, namely, of the constitutive properties of the system.
It may be admitted that on repeated measurement of those conjugate intensive functions of state, they are found to have slightly different values from time to time. Such variability is regarded as due to internal fluctuations. The different measured values average to their nominal values.
If the system is truly macroscopic as postulated by classical thermodynamics, then the fluctuations are too small to detect macroscopically. This is called the thermodynamic limit. In effect, the molecular nature of matter and the quantal nature of momentum transfer have vanished from sight, too small to see. According to Buchdahl: "... there is no place within the strictly phenomenological theory for the idea of fluctuations about equilibrium (see, however, Section 76)."
If the system is repeatedly subdivided, eventually a system is produced that is small enough to exhibit obvious fluctuations. This is a mesoscopic level of investigation. The fluctuations are then directly dependent on the natures of the various walls of the system. The precise choice of independent state variables is then important. At this stage, statistical features of the laws of thermodynamics become apparent.
If the mesoscopic system is further repeatedly divided, eventually a microscopic system is produced. Then the molecular character of matter and the quantal nature of momentum transfer become important in the processes of fluctuation. One has left the realm of classical or macroscopic thermodynamics, and one needs quantum statistical mechanics. The fluctuations can become relatively dominant, and questions of measurement become important.
The statement that 'the system is its own internal thermodynamic equilibrium' may be taken to mean that 'indefinitely many such measurements have been taken from time to time, with no trend in time in the various measured values'. Thus the statement, that 'a system is in its own internal thermodynamic equilibrium, with stated nominal values of its functions of state conjugate to its specifying state variables', is far far more informative than a statement that 'a set of single simultaneous measurements of those functions of state have those same values'. This is because the single measurements might have been made during a slight fluctuation, away from another set of nominal values of those conjugate intensive functions of state, that is due to unknown and different constitutive properties. A single measurement cannot tell whether that might be so, unless there is also knowledge of the nominal values that belong to the equilibrium state.
Thermal equilibrium
An explicit distinction between 'thermal equilibrium' and 'thermodynamic equilibrium' is made by B. C. Eu. He considers two systems in thermal contact, one a thermometer, the other a system in which there are several occurring irreversible processes, entailing non-zero fluxes; the two systems are separated by a wall permeable only to heat. He considers the case in which, over the time scale of interest, it happens that both the thermometer reading and the irreversible processes are steady. Then there is thermal equilibrium without thermodynamic equilibrium. Eu proposes consequently that the zeroth law of thermodynamics can be considered to apply even when thermodynamic equilibrium is not present; also he proposes that if changes are occurring so fast that a steady temperature cannot be defined, then "it is no longer possible to describe the process by means of a thermodynamic formalism. In other words, thermodynamics has no meaning for such a process." This illustrates the importance for thermodynamics of the concept of temperature.
Thermal equilibrium is achieved when two systems in thermal contact with each other cease to have a net exchange of energy. It follows that if two systems are in thermal equilibrium, then their temperatures are the same.
Thermal equilibrium occurs when a system's macroscopic thermal observables have ceased to change with time. For example, an ideal gas whose distribution function has stabilised to a specific Maxwell–Boltzmann distribution would be in thermal equilibrium. This outcome allows a single temperature and pressure to be attributed to the whole system. For an isolated body, it is quite possible for mechanical equilibrium to be reached before thermal equilibrium is reached, but eventually, all aspects of equilibrium, including thermal equilibrium, are necessary for thermodynamic equilibrium.
Non-equilibrium
A system's internal state of thermodynamic equilibrium should be distinguished from a "stationary state" in which thermodynamic parameters are unchanging in time but the system is not isolated, so that there are, into and out of the system, non-zero macroscopic fluxes which are constant in time.
Non-equilibrium thermodynamics is a branch of thermodynamics that deals with systems that are not in thermodynamic equilibrium. Most systems found in nature are not in thermodynamic equilibrium because they are changing or can be triggered to change over time, and are continuously and discontinuously subject to flux of matter and energy to and from other systems. The thermodynamic study of non-equilibrium systems requires more general concepts than are dealt with by equilibrium thermodynamics. Many natural systems still today remain beyond the scope of currently known macroscopic thermodynamic methods.
Laws governing systems which are far from equilibrium are also debatable. One of the guiding principles for these systems is the maximum entropy production principle. It states that a non-equilibrium system evolves such as to maximize its entropy production.
| Physical sciences | Thermodynamics | Physics |
265835 | https://en.wikipedia.org/wiki/Francolin | Francolin | Francolins are birds in the tribe Gallini that traditionally have been placed in the genus Francolinus, but now commonly are divided into multiple genera.
As previously defined, they were paraphyletic as the genus Pternistis, which was previously included in Francolinus, is more closely related to Old World quails than it is to the other francolins. Beginning in 2004, various ornithologists have recommended that it would be clearer to use "spurfowl" for all members of the genus Pternistis and restrict the use of "francolin" to the other species presently or formerly classified in Francolinus. When Pternistis is excluded, the francolins form a monophyletic clade that is a sister group to a clade comprising the junglefowl (Gallus) and the bamboo partridges (Bambusicola); together, these clades compose the tribe Gallini.
Although formerly classified in the partridge subfamily Perdicinae, this classification is no longer supported, and they are now classified in the subfamily Pavoninae.
Francolins are terrestrial (though not flightless) birds that feed on insects, vegetable matter and seeds. Most of the members have a hooked upper beak, well-suited for digging at the bases of grass tussocks and rootballs. They have wide tails with fourteen rectrix feathers. Most species exhibit spurs on the tarsi.
Distribution
Of the approximately 17 extant species, the natural range of five (composing the genus Francolinus and Ortygornis) are restricted to Asia, while the remaining genera are restricted to Africa. Several species have been introduced to other parts of the world, notably Hawaii.
Twelve of the species which occur in Africa are found in the subcontinental region of southern Africa; of these, seven occur in varying proportions within the political boundaries of Namibia and Zambia. Six southern African francolins are considered endemic to the subcontinent, of which three are found in Namibia and Zambia (the Hartlaub's, red-billed and Orange River francolins).
The Cape spurfowl, endemic to the Cape Province of South Africa, occurs marginally in southern Namibia and southwestern Zambia. A fossil francolin, Francolinus capeki, has been described from Late Pliocene deposits of Hungary; the contemporary fossil galliforms "Francolinus" minor and "F." subfrancolinus are now placed in Palaeocryptonyx.
Taxonomy
Until the early 1990s, major authorities placed all francolins in the genus Francolinus. In 1992 it was suggested that this treatment was problematic, and the francolins should be split into four genera: Francolinus for the Asian species, and the African species divided into Peliperdix, Scleroptila and Pternistis. The crested francolin and Nahan's francolin were considered possibly quite distinct, but still maintained in Peliperdix and Pternistis respectively. Based on further evidence, the crested francolin was moved to the monotypic genus Dendroperdix in 1998, and the Nahan's francolin was moved to Ptilopachus in 2006. Though some still maintain all these in Francolinus, the split into multiple genera is becoming more widespread. In 2021, two species in Francolinus (the grey and swamp francolins) along with the crested francolin were moved into the genus Ortygornis, while three species from Peliperdix (the coqui, white-throated, and Schlegel's francolins) were moved into the new genus Campocolinus. Pternistis was moved to the tribe Coturnicini and, as with Nahan's "francolin", is no longer considered a francolin.
When split, the English name "francolin" is generally restricted to the members of the genera Francolinus, Ortygornis, Campocolinus, Peliperdix and Scleroptila, while the name "spurfowl" is used for Pternistis ("spurfowl" is also used for Galloperdix of the Indian subcontinent). As the Nahan's "francolin" is related to the stone partridge rather than the true francolins and spurfowl, its name is sometimes modified to Nahan's partridge.
In addition to the major changes proposed at genus level, the species level taxonomy among several francolins/spurfowl is disputed. For example, the distribution of the Orange River francolin (Scleroptila levaillantoides) is highly disjunct, leading some authorities to split the northern taxa (from Kenya and northwards) into a separate species, the acacia/Archer's francolin (S. gutturalis, with subspecies lorti), while maintaining the southern taxa (from Angola and southwards) in the Orange River francolin. Most authorities treat the Elgon francolin (S. psilolaema elgonensis) as a subspecies of the moorland francolin, but others have suggested it is a species (S. elgonensis), a subspecies of the Shelley's francolin, or even a hybrid between the moorland and red-winged francolins.
Species
Peliperdix:
Latham's francolin, Peliperdix lathami
Ortygornis:
Crested francolin, Ortygornis sephaena
Grey francolin, Ortygornis pondicerianus
Swamp francolin, Ortygornis gularis
Francolinus:
Black francolin, Francolinus francolinus
Painted francolin, Francolinus pictus
Chinese francolin, Francolinus pintadeanus
Campocolinus:
Coqui francolin, Campocolinus coqui
White-throated francolin, Campocolinus albogularis
Schlegel's francolin, Campocolinus schlegelii
Scleroptila:
Ring-necked francolin, Scleroptila streptophora
Red-winged francolin, Scleroptila levaillantii
Finsch's francolin, Scleroptila finschi
Moorland francolin, Scleroptila psilolaema
Grey-winged francolin, Scleroptila afra
Orange River francolin, Scleroptila gutturalis
Shelley's francolin, Scleroptila shelleyi
With a paraphyletic classification, the genus Pternistis in the tribe Coturnicini is also considered a francolin; due to the resulting paraphyly, this classification is no longer supported and has been recommended against.
| Biology and health sciences | Galliformes | Animals |
266319 | https://en.wikipedia.org/wiki/Gnatcatcher | Gnatcatcher | The gnatcatchers are a family of small passerine birds called Polioptilidae. The 22 species occur in North and South America (except for the far south and the high Andean regions). Most species of this mainly tropical and subtropical group are resident, but the blue-grey gnatcatcher of the United States and southern Canada migrates south in winter. They are close relatives of the wrens.
Description
These dainty birds are intermediate between Old World warblers and wrens in their structure and habits, moving restlessly through foliage seeking insects. The gnatcatchers are mainly soft bluish grey in color, and have the typical insectivore's long sharp bill. Many species have distinctive black head patterns (especially males) and long, regularly cocked, black-and-white tails. The skulking gnatwrens are browner, more thickset, and with proportionally shorter tails and longer bills.
Distribution and habitat
They are distributed from North to South America, with the exception of the far south and high Andean regions. Gnatwrens typically occur in the undergrowth of dense, often humid, forest, while gnatcatchers, depending on the species involved, occur in anything from dry scrubby habitats (e.g. the California gnatcatcher) to the canopy of humid Amazonian forest (e.g. the Guianan gnatcatcher). The North American species nest in bushes or trees, but the breeding behavior of several of the Neotropical species is essentially unknown.
Taxonomy and systematics
A species new to science, the critically endangered Iquitos gnatcatcher Polioptila clementsi, was first described in 2005. This species is a member of the Guianan gnatcatcher Polioptila guianensis complex, which recently has been proposed split into three species (four with the Iquitos gnatcatcher), but not all authorities have accepted this (e.g. SACC). Furthermore, other groups should possibly be split, notably the tropical gnatcatcher Polioptila plumbea and masked gnatcatcher Polioptila dumicola complexes, but at present scientific papers on these matters are lacking.
The family contains 22 species divided into 3 genera:
Ramphocaenus (2 species) – gnatwrens
Microbates (2 species) – gnatwrens
Polioptila (18 species) – gnatcatchers
California Gnatcatcher
The California gnatcatcher (Polioptila californica) is a small, non-migratory songbird that inhabits the coastal regions of southern California. It belongs to the family Polioptilidae and is best known for its distinctive behavior and appearance. This species is characterized by a small, slender body, typically around 10 cm in length, with a long tail, dark gray or bluish-gray plumage, and a faint black line running through its eye. The California gnatcatcher is an important species in the coastal sage scrub ecosystem, which it relies on for both food and shelter.
Natural habitat
The natural habitat of the California gnatcatcher is primarily the coastal sage scrub ecosystem. This type of habitat is found in areas with a Mediterranean climate—mild, wet winters and hot, dry summers. The gnatcatcher favors areas with low shrubs and a mix of grasses and wildflowers, where it can easily find insects. It also lives in chaparral and mixed woodlands but is most commonly found below 1,000 meters in elevation. The coastal sage scrub once stretched across much of southern California, from Santa Barbara to San Diego, but has been drastically reduced by urban sprawl and agricultural development.
Food
The California gnatcatcher feeds primarily on insects, such as small flies, beetles, ants, and other invertebrates. It uses its sharp beak to catch these insects in the dense vegetation of its habitat. The gnatcatcher is particularly skilled at foraging in the low shrubbery and in the underbrush, where its small size allows it to maneuver with agility. The species is an insectivore, and its diet is essential for providing the energy it needs for survival and reproduction. It also consumes some seeds and berries, but insects make up the bulk of its diet.
Population
Over the years, the population of the California gnatcatcher has fluctuated dramatically. In the early 1990s, the population was estimated to be only 1,000 to 3,000 pairs due to habitat destruction, primarily from urban expansion. As southern California grew, vast areas of coastal sage scrub were converted into housing developments, agriculture, and infrastructure, which resulted in the loss of the gnatcatcher's habitat. The species was listed as threatened under the U.S. Endangered Species Act in 1993, prompting conservation efforts to protect its remaining habitats. Today, the population is estimated to be around 20,000 to 40,000 individuals. thanks to restoration and preservation projects.
Status
The California gnatcatcher is classified as a threatened species, primarily due to habitat loss and fragmentation. Urbanization, agriculture, and climate change have significantly reduced the available coastal sage scrub habitat. Habitat fragmentation isolates populations, making it difficult for gnatcatchers to find mates and reproduce, which also reduces genetic diversity. While some recovery has been made, the species is still vulnerable, and its long-term survival depends on continued conservation efforts. Protected areas and habitat restoration have been critical in stabilizing the population in recent years.
Behaviorally, California gnatcatchers are highly territorial, especially during the breeding season. They use a combination of vocalizations and physical displays to defend their territories from other birds. These birds are also known for their distinctive foraging behavior. They often move quickly through shrubs, flicking their tails and catching insects in mid-air or from leaves. The gnatcatcher is an agile and persistent forager, spending much of its day hunting for food to sustain itself. The species is also known for being relatively solitary outside of the breeding season, with pairs coming together only for nesting.
Predators
As small birds, California gnatcatchers face several natural predators. They are vulnerable to a variety of avian predators, including hawks, kestrels, and ravens, which may prey upon them or their eggs. Additionally, snakes and raccoons are known to threaten their nests. These predators often target the eggs or fledglings, making early life particularly dangerous for the gnatcatcher. In order to avoid predators, gnatcatchers rely on their excellent camouflage and the dense vegetation of the coastal sage scrub to hide from potential threats.
California gnatcatchers are also particularly vulnerable to wildfires, which have become more frequent and intense due to climate change. Wildfires can destroy vast swaths of their habitat, leaving them without food and shelter. These fires can decimate entire populations and disrupt their breeding cycles. To ensure their survival, California gnatcatchers depend on the preservation of their habitat and the continuation of conservation efforts that focus on restoring degraded areas and protecting the species' remaining habitat from further destruction.
| Biology and health sciences | Passerida | Animals |
266443 | https://en.wikipedia.org/wiki/Metalworking | Metalworking | Metalworking is the process of shaping and reshaping metals in order to create useful objects, parts, assemblies, and large scale structures. As a term, it covers a wide and diverse range of processes, skills, and tools for producing objects on every scale: from huge ships, buildings, and bridges, down to precise engine parts and delicate jewelry.
The historical roots of metalworking predate recorded history; its use spans cultures, civilizations and millennia. It has evolved from shaping soft, native metals like gold with simple hand tools, through the smelting of ores and hot forging of harder metals like iron, up to and including highly technical modern processes such as machining and welding. It has been used as an industry, a driver of trade, individual hobbies, and in the creation of art; it can be regarded as both a science and a craft.
Modern metalworking processes, though diverse and specialized, can be categorized into one of three broad areas known as forming, cutting, or joining processes. Modern metalworking workshops, typically known as machine shops, hold a wide variety of specialized or general-use machine tools capable of creating highly precise, useful products. Many simpler metalworking techniques, such as blacksmithing, are no longer economically competitive on a large scale in developed countries; some of them are still in use in less developed countries, for artisanal or hobby work, or for historical reenactment.
Prehistory
The oldest archaeological evidence of copper mining and working was the discovery of a copper pendant in northern Iraq from 8,700 BCE. The earliest substantiated and dated evidence of metalworking in the Americas was the processing of copper in Wisconsin, near Lake Michigan. Copper was hammered until it became brittle, then heated so it could be worked further. In America, this technology is dated to about 4000–5000 BCE. The oldest gold artifacts in the world come from the Bulgarian Varna Necropolis and date from 4450 BCE.
Not all metal required fire to obtain it or work it. Isaac Asimov speculated that gold was the "first metal". His reasoning being, that, by its chemistry, it is found in nature as nuggets of pure gold. In other words, gold, as rare as it is, is sometimes found in nature as a native metal. Some metals can also be found in meteors. Almost all other metals are found in ores, a mineral-bearing rock, that require heat or some other process to liberate the metal. Another feature of gold is that it is workable as it is found, meaning that no technology beyond a stone hammer and anvil is needed to work the metal. This is a result of gold's properties of malleability and ductility. The earliest tools were stone, bone, wood, and sinew, all of which sufficed to work gold.
At some unknown time, the process of liberating metals from rock by heat became known, and rocks rich in copper, tin, and lead came into demand. These ores were mined wherever they were recognized. Remnants of such ancient mines have been found all over Southwestern Asia. Metalworking was being carried out by the South Asian inhabitants of Mehrgarh between 7000 and 3300 BCE. The end of the beginning of metalworking occurs sometime around 6000 BCE when copper smelting became common in Southwestern Asia.
Ancient civilisations knew of seven metals. Here they are arranged in order of their oxidation potential (in volts):
Iron +0.44 V,
Tin +0.14 V
Lead +0.13 V
Copper −0.34 V
Mercury −0.79 V
Silver −0.80 V
Gold −1.50 V.
The oxidation potential is important because it is one indicator of how tightly bound to the ore the metal is likely to be. As can be seen, iron is significantly higher than the other six metals while gold is dramatically lower than the six above it. Gold's low oxidation is one of the main reasons that gold is found in nuggets. These nuggets are relatively pure gold and are workable as they are found.
Copper ore, being relatively abundant, and tin ore became the next important substances in the story of metalworking. Using heat to smelt copper from ore, a great deal of copper was produced. It was used for both jewelry and simple tools. However, copper by itself was too soft for tools requiring edges and stiffness. At some point tin was added into the molten copper and bronze was developed thereby. Bronze is an alloy of copper and tin. Bronze was an important advance because it had the edge-durability and stiffness that pure copper lacked. Until the advent of iron, bronze was the most advanced metal for tools and weapons in common use (see Bronze Age for more detail).
Outside Southwestern Asia, these same advances and materials were being discovered and used around the world. People in China and Great Britain began using bronze with little time being devoted to copper. Japanese began the use of bronze and iron almost simultaneously. In the Americas it was different. Although the peoples of the Americas knew of metals, it was not until the European colonisation that metalworking for tools and weapons became common. Jewelry and art were the principal uses of metals in the Americas prior to European influence.
About 2700 BCE, production of bronze was common in locales where the necessary materials could be assembled for smelting, heating, and working the metal. Iron was beginning to be smelted and began its emergence as an important metal for tools and weapons. The period that followed became known as the Iron Age.
History
By the historical periods of the Pharaohs in Egypt, the Vedic Kings in India, the Tribes of Israel, and the Maya civilization in North America, among other ancient populations, precious metals began to have value attached to them. In some cases rules for ownership, distribution, and trade were created, enforced, and agreed upon by the respective peoples. By the above periods metalworkers were very skilled at creating objects of adornment, religious artifacts, and trade instruments of precious metals (non-ferrous), as well as weaponry usually of ferrous metals and/or alloys. These skills were well executed. The techniques were practiced by artisans, blacksmiths, atharvavedic practitioners, alchemists, and other categories of metalworkers around the globe. For example, the granulation technique was employed by numerous ancient cultures before the historic record shows people traveled to far regions to share this process. Metalsmiths today still use this and many other ancient techniques.
As time progressed, metal objects became more common, and ever more complex. The need to further acquire and work metals grew in importance. Skills related to extracting metal ores from the earth began to evolve, and metalsmiths became more knowledgeable. Metalsmiths became important members of society. Fates and economies of entire civilizations were greatly affected by the availability of metals and metalsmiths. The metalworker depends on the extraction of precious metals to make jewelry, build more efficient electronics, and for industrial and technological applications from construction to shipping containers to rail, and air transport. Without metals, goods and services would cease to move around the globe on the scale we know today.
General processes
Metalworking generally is divided into three categories: forming, cutting, and joining. Most metal cutting is done by high speed steel tools or carbide tools. Each of these categories contains various processes.
Prior to most operations, the metal must be marked out and/or measured, depending on the desired finished product.
Marking out (also known as layout) is the process of transferring a design or pattern to a workpiece and is the first step in the handcraft of metalworking. It is performed in many industries or hobbies, although in industry, the repetition eliminates the need to mark out every individual piece. In the metal trades area, marking out consists of transferring the engineer's plan to the workpiece in preparation for the next step, machining or manufacture.
Calipers are hand tools designed to precisely measure the distance between two points. Most calipers have two sets of flat, parallel edges used for inner or outer diameter measurements. These calipers can be accurate to within one-thousandth of an inch (25.4 μm). Different types of calipers have different mechanisms for displaying the distance measured. Where larger objects need to be measured with less precision, a tape measure is often used.
Casting
Casting achieves a specific form by pouring molten metal into a mold and allowing it to cool, with no mechanical force. Forms of casting include:
Investment casting (called lost wax casting in art)
Centrifugal casting
Die casting
Sand casting
Shell casting
Spin casting
Forming processes
These forming processes modify metal or workpiece by deforming the object, that is, without removing any material. Forming is done with a system of mechanical forces and, especially for bulk metal forming, with heat.
Bulk forming processes
Plastic deformation involves using heat or pressure to make a workpiece more conductive to mechanical force. Historically, this and casting were done by blacksmiths, though today the process has been industrialized. In bulk metal forming, the workpiece is generally heated up.
Cold sizing
Extrusion
Drawing
Forging
Powder metallurgy
Friction drilling
Rolling
Burnishing
Sheet (and tube) forming processes
These types of forming process involve the application of mechanical force at room temperature. However, some recent developments involve the heating of dies and/or parts. Advancements in automated metalworking technology have made progressive die stamping possible which is a method that can encompass punching, coining, bending and several other ways below that modify metal at less cost while resulting in less scrap.
Bending
Coining
Decambering
Deep drawing (DD)
Foldforming
Hydroforming (HF)
Hot metal gas forming
Hot press hardening
Incremental forming (IF)
Spinning, Shear forming or Flowforming
Planishing
Raising
Roll forming
Roll bending
Repoussé and chasing
Rubber pad forming
Shearing
Stamping
Superplastic forming (SPF)
Wheeling using an English wheel (wheeling machine)
Cutting processes
Cutting is a collection of processes wherein material is brought to a specified geometry by removing excess material using various kinds of tooling to leave a finished part that meets specifications. The net result of cutting is two products, the waste or excess material, and the finished part. In woodworking, the waste would be sawdust and excess wood. In cutting metals the waste is chips or swarf and excess metal.
Cutting processes fall into one of three major categories:
Chip producing processes most commonly known as machining
Burning, a set of processes wherein the metal is cut by oxidizing a kerf to separate pieces of metal
Miscellaneous specialty process, not falling easily into either of the above categories
Drilling a hole in a metal part is the most common example of a chip producing process. Using an oxy-fuel cutting torch to separate a plate of steel into smaller pieces is an example of burning. Chemical milling is an example of a specialty process that removes excess material by the use of etching chemicals and masking chemicals.
There are many technologies available to cut metal, including:
Manual technologies: saw, chisel, shear or snips
Machine technologies: turning, milling, drilling, grinding, sawing
Welding/burning technologies: burning by laser, oxy-fuel burning, and plasma
Erosion technologies: by water jet, electric discharge, or abrasive flow machining.
Chemical technologies: Photochemical machining
Cutting fluid or coolant is used where there is significant friction and heat at the cutting interface between a cutter such as a drill or an end mill and the workpiece. Coolant is generally introduced by a spray across the face of the tool and workpiece to decrease friction and temperature at the cutting tool/workpiece interface to prevent excessive tool wear. In practice there are many methods of delivering coolant.
Health effects
The use of an angle grinder in cutting is not preferred as large amounts of harmful sparks and fumes (and particulates) are generated when compared with using reciprocating saw or band saw. Angle grinders produce sparks when cutting ferrous metals. They also produce shards cutting other materials.
Milling
Milling is the complex shaping of metal or other materials by removing material to form the final shape. It is generally done on a milling machine, a power-driven machine that in its basic form consists of a milling cutter that rotates about the spindle axis (like a drill), and a worktable that can move in multiple directions (usually two dimensions [x and y axis] relative to the workpiece). The spindle usually moves in the z axis. It is possible to raise the table (where the workpiece rests). Milling machines may be operated manually or under computer numerical control (CNC), and can perform a vast number of complex operations, such as slot cutting, planing, drilling and threading, rabbeting, routing, etc. Two common types of mills are the horizontal mill and vertical mill.
The pieces produced are usually complex 3D objects that are converted into x, y, and z coordinates that are then fed into the CNC machine and allow it to complete the tasks required. The milling machine can produce most parts in 3D, but some require the objects to be rotated around the x, y, or z coordinate axis (depending on the need). Tolerances come in a variety of standards, depending on the locale. In countries still using the imperial system, this is usually in the thousandths of an inch (unit known as thou), depending on the specific machine. In many other European countries, standards following the ISO are used instead.
In order to keep both the bit and material cool, a high temperature coolant is used. In most cases the coolant is sprayed from a hose directly onto the bit and material. This coolant can either be machine or user controlled, depending on the machine.
Materials that can be milled range from aluminum to stainless steel and almost everything in between. Each material requires a different speed on the milling tool and varies in the amount of material that can be removed in one pass of the tool. Harder materials are usually milled at slower speeds with small amounts of material removed. Softer materials vary, but usually are milled with a high bit speed.
The use of a milling machine adds costs that are factored into the manufacturing process. Each time the machine is used coolant is also used, which must be periodically added in order to prevent breaking bits. A milling bit must also be changed as needed in order to prevent damage to the material. Time is the biggest factor for costs. Complex parts can require hours to complete, while very simple parts take only minutes. This in turn varies the production time as well, as each part will require different amounts of time.
Safety is key with these machines. The bits are traveling at high speeds and removing pieces of usually scalding hot metal. The advantage of having a CNC milling machine is that it protects the machine operator.
Turning
Turning is a metal cutting process for producing a cylindrical surface with a single point tool. The workpiece is rotated on a spindle and the cutting tool is fed into it radially, axially or both. Producing surfaces perpendicular to the workpiece axis is called facing. Producing surfaces using both radial and axial feeds is called profiling.
A lathe is a machine tool which spins a block or cylinder of material so that when abrasive, cutting, or deformation tools are applied to the workpiece, it can be shaped to produce an object which has rotational symmetry about an axis of rotation. Examples of objects that can be produced on a lathe include candlestick holders, crankshafts, camshafts, and bearing mounts.
Lathes have four main components: the bed, the headstock, the carriage, and the tailstock. The bed is a precise & very strong base which all of the other components rest upon for alignment. The headstock's spindle secures the workpiece with a chuck, whose jaws (usually three or four) are tightened around the piece. The spindle rotates at high speed, providing the energy to cut the material. While historically lathes were powered by belts from a line shaft, modern examples uses electric motors. The workpiece extends out of the spindle along the axis of rotation above the flat bed. The carriage is a platform that can be moved, precisely and independently parallel and perpendicular to the axis of rotation. A hardened cutting tool is held at the desired height (usually the middle of the workpiece) by the toolpost. The carriage is then moved around the rotating workpiece, and the cutting tool gradually removes material from the workpiece. The tailstock can be slid along the axis of rotation and then locked in place as necessary. It may hold centers to further secure the workpiece, or cutting tools driven into the end of the workpiece.
Other operations that can be performed with a single point tool on a lathe are:
Chamfering: Cutting an angle on the corner of a cylinder.
Parting: The tool is fed radially into the workpiece to cut off the end of a part.
Threading: A tool is fed along and across the outside or inside surface of rotating parts to produce external or internal threads.
Boring: A single-point tool is fed linearly and parallel to the axis of rotation to create a round hole.
Drilling: Feeding the drill into the workpiece axially.
Knurling: Uses a tool to produce a rough surface texture on the work piece. Frequently used to allow grip by hand on a metal part.
Modern computer numerical control (CNC) lathes and (CNC) machining centres can do secondary operations like milling by using driven tools. When driven tools are used the work piece stops rotating and the driven tool executes the machining operation with a rotating cutting tool. The CNC machines use x, y, and z coordinates in order to control the turning tools and produce the product. Most modern day CNC lathes are able to produce most turned objects in 3D.
Nearly all types of metal can be turned, although more time & specialist cutting tools are needed for harder workpieces.
Threading
There are many threading processes including: cutting threads with a tap or die, thread milling, single-point thread cutting, thread rolling, cold root rolling and forming, and thread grinding. A tap is used to cut a female thread on the inside surface of a pre-drilled hole, while a die cuts a male thread on a preformed cylindrical rod.
Grinding
Grinding uses an abrasive process to remove material from the workpiece. A grinding machine is a machine tool used for producing very fine finishes, making very light cuts, or high precision forms using an abrasive wheel as the cutting device. This wheel can be made up of various sizes and types of stones, diamonds or inorganic materials.
The simplest grinder is a bench grinder or a hand-held angle grinder, for deburring parts or cutting metal with a zip-disc.
Grinders have increased in size and complexity with advances in time and technology. From the old days of a manual toolroom grinder sharpening endmills for a production shop, to today's 30000 RPM CNC auto-loading manufacturing cell producing jet turbines, grinding processes vary greatly.
Grinders need to be very rigid machines to produce the required finish. Some grinders are even used to produce glass scales for positioning CNC machine axis. The common rule is the machines used to produce scales be 10 times more accurate than the machines the parts are produced for.
In the past grinders were used for finishing operations only because of limitations of tooling. Modern grinding wheel materials and the use of industrial diamonds or other man-made coatings (cubic boron nitride) on wheel forms have allowed grinders to achieve excellent results in production environments instead of being relegated to the back of the shop.
Modern technology has advanced grinding operations to include CNC controls, high material removal rates with high precision, lending itself well to aerospace applications and high volume production runs of precision components.
Filing
Filing is combination of grinding and saw tooth cutting using a file. Prior to the development of modern machining equipment it provided a relatively accurate means for the production of small parts, especially those with flat surfaces. The skilled use of a file allowed a machinist to work to fine tolerances and was the hallmark of the craft. Today filing is rarely used as a production technique in industry, though it remains as a common method of deburring.
Other
Broaching is a machining operation used to cut keyways into shafts. Electron beam machining (EBM) is a machining process where high-velocity electrons are directed toward a work piece, creating heat and vaporizing the material. Ultrasonic machining uses ultrasonic vibrations to machine very hard or brittle materials.
Joining processes
Welding
Welding is a fabrication process that joins materials, usually metals or thermoplastics, by causing coalescence. This is often done by melting the workpieces and adding a filler material to form a pool of molten material that cools to become a strong joint, but sometimes pressure is used in conjunction with heat, or by itself, to produce the weld.
Many different energy sources can be used for welding, including a gas flame, an electric arc, a laser, an electron beam, friction, and ultrasound. While often an industrial process, welding can be done in many different environments, including open air, underwater and in space. Regardless of location, however, welding remains dangerous, and precautions must be taken to avoid burns, electric shock, poisonous fumes, and overexposure to ultraviolet light.
Brazing
Brazing is a joining process in which a filler metal is melted and drawn into a capillary formed by the assembly of two or more work pieces. The filler metal reacts metallurgically with the workpieces and solidifies in the capillary, forming a strong joint. Unlike welding, the work piece is not melted. Brazing is similar to soldering, but occurs at temperatures in excess of . Brazing has the advantage of producing less thermal stresses than welding, and brazed assemblies tend to be more ductile than weldments because alloying elements can not segregate and precipitate.
Brazing techniques include, flame brazing, resistance brazing, furnace brazing, diffusion brazing, inductive brazing and vacuum brazing.
Soldering
Soldering is a joining process that occurs at temperatures below . It is similar to brazing in the way that a filler is melted and drawn into a capillary to form a joint, although at a lower temperature. Because of this lower temperature and different alloys used as fillers, the metallurgical reaction between filler and work piece is minimal, resulting in a weaker joint.
Riveting
Riveting is one of the most ancient metalwork joining processes. Its use declined markedly during the second half of the 20th century, but it still retains important uses in industry and construction, and in artisan crafts such as jewellery, medieval armouring and metal couture in the early 21st century. The earlier use of rivets is being superseded by improvements in welding and component fabrication techniques.
A rivet is essentially a two-headed and unthreaded bolt which holds two other pieces of metal together. Holes are drilled or punched through the two pieces of metal to be joined. The holes being aligned, a rivet is passed through the holes and permanent heads are formed onto the ends of the rivet utilizing hammers and forming dies (by either cold working or hot working).
Rivets are commonly purchased with one head already formed.
When it is necessary to remove rivets, one of the rivet's heads is sheared off with a cold chisel. The rivet is then driven out with a hammer and punch.
Mechanical fixings
This includes screws, as well as bolts. This is often used as it requires relatively little specialist equipment, and are therefore often used in flat-pack furniture. It can also be used when a metal is joined to another material (such as wood) or a particular metal does not weld well (such as aluminum). This can be done to directly join metals, or with an intermediate material such as nylon. While often weaker than other methods such as welding or brazing, the metal can easily be removed and therefore reused or recycled. It can also be done in conjunction with an epoxy or glue, reverting its ecological benefits.
Associated processes
While these processes are not primary metalworking processes, they are often performed before or after metalworking processes.
Heat treatment
Metals can be heat treated to alter the properties of strength, ductility, toughness, hardness or resistance to corrosion. Common heat treatment processes include annealing, precipitation hardening, quenching, and tempering:
annealing softens the metal by allowing recovery of cold work and grain growth.
quenching can be used to harden alloy steels, or in precipitation hardenable alloys, to trap dissolved solute atoms in solution.
tempering will cause the dissolved alloying elements to precipitate, or in the case of quenched steels, improve impact strength and ductile properties.
Often, mechanical and thermal treatments are combined in what is known as thermo-mechanical treatments for better properties and more efficient processing of materials. These processes are common to high alloy special steels, super alloys and titanium alloys.
Plating
Electroplating is a common surface-treatment technique. It involves bonding a thin layer of another metal such as gold, silver, chromium or zinc to the surface of the product by hydrolysis. It is used to reduce corrosion, create abrasion resistance and improve the product's aesthetic appearance. Plating can even change the properties of the original part including conductivity, heat dissipation or structural integrity. There are four main electroplating methods to ensure proper coating and cost effectiveness per product: mass plating, rack plating, continuous plating and line plating.
Thermal spraying
Thermal spraying techniques are another popular finishing option, and often have better high temperature properties than electroplated coatings due to the thicker coating. The four main thermal spray processes include electric wire arc spray, flame (oxy acetylene combustion) spray, plasma spray and high velocity oxy fuel (HVOF) spray.
| Technology | Metallurgy | null |
266466 | https://en.wikipedia.org/wiki/Octet%20rule | Octet rule | The octet rule is a chemical rule of thumb that reflects the theory that main-group elements tend to bond in such a way that each atom has eight electrons in its valence shell, giving it the same electronic configuration as a noble gas. The rule is especially applicable to carbon, nitrogen, oxygen, and the halogens; although more generally the rule is applicable for the s-block and p-block of the periodic table. Other rules exist for other elements, such as the duplet rule for hydrogen and helium, and the 18-electron rule for transition metals.
The valence electrons can be counted using a Lewis electron dot diagram as shown at the right for carbon dioxide. The electrons shared by the two atoms in a covalent bond are counted twice, once for each atom. In carbon dioxide each oxygen shares four electrons with the central carbon, two (shown in red) from the oxygen itself and two (shown in black) from the carbon. All four of these electrons are counted in both the carbon octet and the oxygen octet, so that both atoms are considered to obey the octet rule.
Example: sodium chloride (NaCl)
Ionic bonding is common between pairs of atoms, where one of the pair is a metal of low electronegativity (such as sodium) and the second a nonmetal of high electronegativity (such as chlorine).
A chlorine atom has seven electrons in its third and outer electron shell, the first and second shells being filled with two and eight electrons respectively. The first electron affinity of chlorine (the energy release when chlorine gains an electron to form Cl−) is 349 kJ per mole of chlorine atoms. Adding a second electron to form a hypothetical Cl2- would require energy, energy that cannot be recovered by the formation of a chemical bond. The result is that chlorine will very often form a compound in which it has eight electrons in its outer shell (a complete octet), as in Cl−.
A sodium atom has a single electron in its outermost electron shell, the first and second shells again being full with two and eight electrons respectively. To remove this outer electron requires only the first ionization energy, which is +495.8 kJ per mole of sodium atoms, a small amount of energy. By contrast, the second electron resides in the deeper second electron shell, and the second ionization energy required for its removal is much larger: +4562 kJ per mole. Thus sodium will, in most cases, form a compound in which it has lost a single electron and have a full outer shell of eight electrons, or octet.
The energy required to transfer an electron from a sodium atom to a chlorine atom (the difference of the 1st ionization energy of sodium and the electron affinity of chlorine) is small: +495.8 − 349 = +147 kJ mol−1. This energy is easily offset by the lattice energy of sodium chloride: −783 kJ mol−1. This completes the explanation of the octet rule in this case.
History
In 1864, the English chemist John Newlands classified the sixty-two known elements into eight groups, based on their physical properties.
In the late 19th century, it was known that coordination compounds (formerly called "molecular compounds") were formed by the combination of atoms or molecules in such a manner that the valencies of the atoms involved apparently became satisfied. In 1893, Alfred Werner showed that the number of atoms or groups associated with a central atom (the "coordination number") is often 4 or 6; other coordination numbers up to a maximum of 8 were known, but less frequent. In 1904, Richard Abegg was one of the first to extend the concept of coordination number to a concept of valence in which he distinguished atoms as electron donors or acceptors, leading to positive and negative valence states that greatly resemble the modern concept of oxidation states. Abegg noted that the difference between the maximum positive and negative valences of an element under his model is frequently eight. In 1916, Gilbert N. Lewis referred to this insight as Abegg's rule and used it to help formulate his cubical atom model and the "rule of eight", which began to distinguish between valence and valence electrons. In 1919, Irving Langmuir refined these concepts further and renamed them the "cubical octet atom" and "octet theory". The "octet theory" evolved into what is now known as the "octet rule".
Walther Kossel and Gilbert N. Lewis saw that noble gases did not have the tendency of taking part in chemical reactions under ordinary conditions. On the basis of this observation, they concluded that atoms of noble gases are stable and on the basis of this conclusion they proposed a theory of valency known as "electronic theory of valency" in 1916:
Explanation in quantum theory
The quantum theory of the atom explains the eight electrons as a closed shell with an s2p6 electron configuration. A closed-shell configuration is one in which low-lying energy levels are full and higher energy levels are empty. For example, the neon atom ground state has a full shell (2s2 2p6) and an empty shell. According to the octet rule, the atoms immediately before and after neon in the periodic table (i.e. C, N, O, F, Na, Mg and Al), tend to attain a similar configuration by gaining, losing, or sharing electrons.
The argon atom has an analogous 3s2 3p6 configuration. There is also an empty 3d level, but it is at considerably higher energy than 3s and 3p (unlike in the hydrogen atom), so that 3s2 3p6 is still considered a closed shell for chemical purposes. The atoms immediately before and after argon tend to attain this configuration in compounds. There are, however, some hypervalent molecules in which the 3d level may play a part in the bonding, although this is controversial (see below).
For helium there is no 1p level according to the quantum theory, so that 1s2 is a closed shell with no p electrons. The atoms before and after helium (H and Li) follow a duet rule and tend to have the same 1s2 configuration as helium.
Exceptions
Many reactive intermediates are unstable and do not obey the octet rule. This includes species such as carbenes, as well as free radicals such as the methyl radical (CH3), which has an unpaired electron in a non-bonding orbital on the carbon atom and no electron of opposite spin in the same orbital. Another example is the radical chlorine monoxide (ClO•) which is involved in ozone depletion. These molecules often react so as to complete their octet. Electron deficient molecules such as boranes also do not obey the octet rule but share delocalized electrons in a manner similar to metallic bonding.
Although stable odd-electron molecules and hypervalent molecules are commonly taught as violating the octet rule, ab initio molecular orbital calculations show that they largely obey the octet rule (see three-electron bonds and hypervalent molecules sections below).
Three-electron bonds
Some stable molecular radicals (e.g. nitric oxide, NO) obtain octet configurations by means of a three-electron bond which contributes one shared and one unshared electron to the octet of each bonded atom. In NO, the octet on each atom consists of two electrons from the three-electron bond, plus four electrons from two two-electron bonds and two electrons from a lone pair of non-bonding electrons on that atom alone. The bond order is 2.5, since each two-electron bond counts as one bond while the three-electron bond has only one shared electron and therefore corresponds to a half-bond.
Dioxygen is sometimes represented as obeying the octet rule with a double bond (O=O) containing two pairs of shared electrons. However the ground state of this molecule is paramagnetic, indicating the presence of unpaired electrons. Pauling proposed that this molecule actually contains two three-electron bonds and one normal covalent (two-electron) bond. The octet on each atom then consists of two electrons from each three-electron bond, plus the two electrons of the covalent bond, plus one lone pair of non-bonding electrons. The bond order is 1+0.5+0.5=2.
Hypervalent molecules
Main-group elements in the third and later rows of the periodic table can form hypercoordinate or hypervalent molecules in which the central main-group atom is bonded to more than four other atoms, such as phosphorus pentafluoride, PF5, and sulfur hexafluoride, SF6. For example, in PF5, if it is supposed that there are five true covalent bonds in which five distinct electron pairs are shared, then the phosphorus would be surrounded by 10 valence electrons in violation of the octet rule. In the early days of quantum mechanics, Pauling proposed that third-row atoms can form five bonds by using one s, three p and one d orbitals, or six bonds by using one s, three p and two d orbitals. To form five bonds, the one s, three p and one d orbitals combine to form five sp3d hybrid orbitals which each share an electron pair with a halogen atom, for a total of 10 shared electrons, two more than the octet rule predicts. Similarly to form six bonds, the six sp3d2 hybrid orbitals form six bonds with 12 shared electrons. In this model the availability of empty d orbitals is used to explain the fact that third-row atoms such as phosphorus and sulfur can form more than four covalent bonds, whereas second-row atoms such as nitrogen and oxygen are strictly limited by the octet rule.
However other models describe the bonding using only s and p orbitals in agreement with the octet rule. A valence bond description of PF5 uses resonance between different PF4+ F− structures, so that each F is bonded by a covalent bond in four structures and an ionic bond in one structure. Each resonance structure has eight valence electrons on P. A molecular orbital theory description considers the highest occupied molecular orbital to be a non-bonding orbital localized on the five fluorine atoms, in addition to four occupied bonding orbitals, so again there are only eight valence electrons on the phosphorus. The validity of the octet rule for hypervalent molecules is further supported by ab initio molecular orbital calculations, which show that the contribution of d functions to the bonding orbitals is small.
Nevertheless, for historical reasons, structures implying more than eight electrons around elements like P, S, Se, or I are still common in textbooks and research articles. In spite of the unimportance of d shell expansion in chemical bonding, this practice allows structures to be shown without using a large number of formal charges or using partial bonds and is recommended by the IUPAC as a convenient formalism in preference to depictions that better reflect the bonding. On the other hand, showing more than eight electrons around Be, B, C, N, O, or F (or more than two around H, He, or Li) is considered an error by most authorities.
Other rules
The octet rule is only applicable to main-group elements. Other elements follow other electron counting rules as their valence electron configurations are different from main-group elements. These other rules are shown below:
The duet rule or duplet rule of the first shell applies to H, He and Li—the noble gas helium has two electrons in its outer shell, which is very stable. (Since there is no 1p subshell, 1s is followed immediately by 2s, and thus shell 1 can only have at most 2 valence electrons). Hydrogen only needs one additional electron to attain this stable configuration, while lithium needs to lose one.
For transition metals, molecules tend to obey the 18-electron rule which corresponds to the utilization of valence d, s and p orbitals to form bonding and non-bonding orbitals. However, unlike the octet rule for main-group elements, transition metals do not strictly obey the 18-electron rule and the valence electron count can vary between 12 and 18.
| Physical sciences | Bonding | Chemistry |
266601 | https://en.wikipedia.org/wiki/Refracting%20telescope | Refracting telescope | A refracting telescope (also called a refractor) is a type of optical telescope that uses a lens as its objective to form an image (also referred to a dioptric telescope). The refracting telescope design was originally used in spyglasses and astronomical telescopes but is also used for long-focus camera lenses. Although large refracting telescopes were very popular in the second half of the 19th century, for most research purposes, the refracting telescope has been superseded by the reflecting telescope, which allows larger apertures. A refractor's magnification is calculated by dividing the focal length of the objective lens by that of the eyepiece.
Refracting telescopes typically have a lens at the front, then a long tube, then an eyepiece or instrumentation at the rear, where the telescope view comes to focus. Originally, telescopes had an objective of one element, but a century later, two and even three element lenses were made.
Refracting telescopes use technology that has often been applied to other optical devices, such as binoculars and zoom lenses/telephoto lens/long-focus lens.
Invention
Refractors were the earliest type of optical telescope. The first record of a refracting telescope appeared in the Netherlands about 1608, when a spectacle maker from Middelburg named Hans Lippershey unsuccessfully tried to patent one. News of the patent spread fast and Galileo Galilei, happening to be in Venice in the month of May 1609, heard of the invention, constructed a version of his own, and applied it to making astronomical discoveries.
Refracting telescope designs
All refracting telescopes use the same principles. The combination of an objective lens 1 and some type of eyepiece 2 is used to gather more light than the human eye is able to collect on its own, focus it 5, and present the viewer with a brighter, clearer, and magnified virtual image 6.
The objective in a refracting telescope refracts or bends light. This refraction causes parallel light rays to converge at a focal point; while those not parallel converge upon a focal plane. The telescope converts a bundle of parallel rays to make an angle α, with the optical axis to a second parallel bundle with angle β. The ratio β/α is called the angular magnification. It equals the ratio between the retinal image sizes obtained with and without the telescope.
Refracting telescopes can come in many different configurations to correct for image orientation and types of aberration. Because the image was formed by the bending of light, or refraction, these telescopes are called refracting telescopes or refractors.
Galilean telescope
The design Galileo Galilei used is commonly called a Galilean telescope. It used a convergent (plano-convex) objective lens and a divergent (plano-concave) eyepiece lens (Galileo, 1610). A Galilean telescope, because the design has no intermediary focus, results in a non-inverted (i.e., upright) image.
Galileo's most powerful telescope, with a total length of just under , magnified objects about 30 times. Galileo had to work with the poor lens technology of the time, and found he had to use aperture stops to reduce the diameter of the objective lens (increase its focal ratio) to limit aberrations, so his telescope produced blurry and distorted images with a narrow field of view. Despite these flaws, the telescope was still good enough for Galileo to explore the sky. He used it to view craters on the Moon, the four largest moons of Jupiter, and the phases of Venus.
Parallel rays of light from a distant object (y) would be brought to a focus in the focal plane of the objective lens (F′ L1 / y′). The (diverging) eyepiece (L2) lens intercepts these rays and renders them parallel once more. Non-parallel rays of light from the object traveling at an angle α1 to the optical axis travel at a larger angle (α2 > α1) after they passed through the eyepiece. This leads to an increase in the apparent angular size and is responsible for the perceived magnification.
The final image (y″) is a virtual image, located at infinity and is the same way up (i.e., non-inverted or upright) as the object.
Keplerian telescope
The Keplerian telescope, invented by Johannes Kepler in 1611, is an improvement on Galileo's design. It uses a convex lens as the eyepiece instead of Galileo's concave one. The advantage of this arrangement is that the rays of light emerging from the eyepiece are converging. This allows for a much wider field of view and greater eye relief, but the image for the viewer is inverted. Considerably higher magnifications can be reached with this design, but, like the Galilean telescope, it still uses simple single element objective lens so needs to have a very high focal ratio to reduce aberrations (Johannes Hevelius built an unwieldy f/225 telescope with a objective and a focal length, and even longer tubeless "aerial telescopes" were constructed). The design also allows for use of a micrometer at the focal plane (to determine the angular size and/or distance between objects observed).
Huygens built an aerial telescope for Royal Society of London with a 19 cm (7.5″) single-element lens.
Achromatic refractors
The next major step in the evolution of refracting telescopes was the invention of the achromatic lens, a lens with multiple elements that helped solve problems with chromatic aberration and allowed shorter focal lengths. It was invented in 1733 by an English barrister named Chester Moore Hall, although it was independently invented and patented by John Dollond around 1758. The design overcame the need for very long focal lengths in refracting telescopes by using an objective made of two pieces of glass with different dispersion, 'crown' and 'flint glass', to reduce chromatic and spherical aberration. Each side of each piece is ground and polished, and then the two pieces are assembled together. Achromatic lenses are corrected to bring two wavelengths (typically red and blue) into focus in the same plane.
Chester More Hall is noted as having made the first twin color corrected lens in 1730.
Dollond achromats were quite popular in the 18th century. A major appeal was they could be made shorter. However, problems with glass making meant that the glass objectives were not made more than about in diameter.
In the late 19th century, the Swiss optician Pierre-Louis Guinand developed a way to make higher quality glass blanks of greater than . He passed this technology to his apprentice Joseph von Fraunhofer, who further developed this technology and also developed the Fraunhofer doublet lens design. The breakthrough in glass making techniques led to the great refractors of the 19th century, that became progressively larger through the decade, eventually reaching over 1 meter by the end of that century before being superseded by silvered-glass reflecting telescopes in astronomy.
Noted lens makers of the 19th century include:
Alvan Clark
Brashear
Chance Brothers
Cauchoix
Fraunhofer
Gautier
Grubb
Henry Brothers
Lerebours
Tulley
Some famous 19th century doublet refractors are the James Lick telescope (91 cm/36 in) and the Greenwich 28 inch refractor (71 cm). An example of an older refractor is the Shuckburgh telescope (dating to the late 1700s). A famous refractor was the "Trophy Telescope", presented at the 1851 Great Exhibition in London. The era of the 'great refractors' in the 19th century saw large achromatic lenses, culminating with the largest achromatic refractor ever built, the Great Paris Exhibition Telescope of 1900.
In the Royal Observatory, Greenwich an 1838 instrument named the Sheepshanks telescope includes an objective by Cauchoix. The Sheepshanks had a wide lens, and was the biggest telescope at Greenwich for about twenty years.
An 1840 report from the Observatory noted of the then-new Sheepshanks telescope with the Cauchoix doublet:In the 1900s a noted optics maker was Zeiss. An example of prime achievements of refractors, over 7 million people have been able to view through the 12-inch Zeiss refractor at Griffith Observatory since its opening in 1935; this is the most people to have viewed through any telescope.
Achromats were popular in astronomy for making star catalogs, and they required less maintenance than metal mirrors. Some famous discoveries using achromats are the planet Neptune and the Moons of Mars.
The long achromats, despite having smaller aperture than the larger reflectors, were often favored for "prestige" observatories. In the late 18th century, every few years, a larger and longer refractor would debut.
For example, the Nice Observatory debuted with refractor, the largest at the time, but was surpassed within only a couple of years.
Apochromatic refractors
Apochromatic refractors have objectives built with special, extra-low dispersion materials. They are designed to bring three wavelengths (typically red, green, and blue) into focus in the same plane. The residual color error (tertiary spectrum) can be an order of magnitude less than that of an achromatic lens. Such telescopes contain elements of fluorite or special, extra-low dispersion (ED) glass in the objective and produce a very crisp image that is virtually free of chromatic aberration. Due to the special materials needed in the fabrication, apochromatic refractors are usually more expensive than telescopes of other types with a comparable aperture.
In the 18th century, Dollond, a popular maker of doublet telescopes, also made a triplet, although they were not really as popular as the two element telescopes.
One of the famous triplet objectives is the Cooke triplet, noted for being able to correct the Seidal aberrations. It is recognized as one of the most important objective designs in the field of photography. The Cooke triplet can correct, with only three elements, for one wavelength, spherical aberration, coma, astigmatism, field curvature, and distortion.
Technical considerations
Refractors suffer from residual chromatic and spherical aberration. This affects shorter focal ratios more than longer ones. An achromatic refractor is likely to show considerable color fringing (generally a purple halo around bright objects); an 16 achromat has much less color fringing.
In very large apertures, there is also a problem of lens sagging, a result of gravity deforming glass. Since a lens can only be held in place by its edge, the center of a large lens sags due to gravity, distorting the images it produces. The largest practical lens size in a refracting telescope is around .
There is a further problem of glass defects, striae or small air bubbles trapped within the glass. In addition, glass is opaque to certain wavelengths, and even visible light is dimmed by reflection and absorption when it crosses the air-glass interfaces and passes through the glass itself. Most of these problems are avoided or diminished in reflecting telescopes, which can be made in far larger apertures and which have all but replaced refractors for astronomical research.
The ISS-WAC on the Voyager 1/2 used a lens, launched into space in the late 1970s, an example of the use of refractors in space.
Applications and achievements
Refracting telescopes were noted for their use in astronomy as well as for terrestrial viewing. Many early discoveries of the Solar System were made with singlet refractors.
The use of refracting telescopic optics are ubiquitous in photography, and are also used in Earth orbit.
One of the more famous applications of the refracting telescope was when Galileo used it to discover the four largest moons of Jupiter in 1609. Furthermore, early refractors were also used several decades later to discover Titan, the largest moon of Saturn, along with three more of Saturn's moons.
In the 19th century, refracting telescopes were used for pioneering work on astrophotography and spectroscopy, and the related instrument, the heliometer, was used to calculate the distance to another star for the first time. Their modest apertures did not lead to as many discoveries and typically so small in aperture that many astronomical objects were simply not observable until the advent of long-exposure photography, by which time the reputation and quirks of reflecting telescopes were beginning to exceed those of the refractors. Despite this, some discoveries include the Moons of Mars, a fifth Moon of Jupiter, and many double star discoveries including Sirius (the Dog star). Refractors were often used for positional astronomy, besides from the other uses in photography and terrestrial viewing.
Singlets
The Galilean moons and many other moons of the solar system, were discovered with single-element objectives and aerial telescopes.
Galileo Galilei's discovered the Galilean satellites of Jupiter in 1610 with a refracting telescope.
The planet Saturn's moon, Titan, was discovered on March 25, 1655, by the Dutch astronomer Christiaan Huygens.
Doublets
In 1861, the brightest star in the night sky, Sirius, was found to have smaller stellar companion using the 18 and half-inch Dearborn refracting telescope.
By the 18th century refractors began to have major competition from reflectors, which could be made quite large and did not normally suffer from the same inherent problem with chromatic aberration. Nevertheless, the astronomical community continued to use doublet refractors of modest aperture in comparison to modern instruments. Noted discoveries include the Moons of Mars and a fifth moon of Jupiter, Amalthea.
Asaph Hall discovered Deimos on 12 August 1877 at about 07:48 UTC and Phobos on 18 August 1877, at the US Naval Observatory in Washington, D.C., at about 09:14 GMT (contemporary sources, using the pre-1925 astronomical convention that began the day at noon, give the time of discovery as 11 August 14:40 and 17 August 16:06 Washington mean time respectively).
The telescope used for the discovery was the refractor (telescope with a lens) then located at Foggy Bottom. In 1893 the lens was remounted and put in a new dome, where it remains into the 21st century.
Jupiter's moon Amalthea was discovered on 9 September 1892, by Edward Emerson Barnard using the refractor telescope at Lick Observatory. It was discovered by direct visual observation with the doublet-lens refractor.
In 1904, one of the discoveries made using Great Refractor of Potsdam (a double telescope with two doublets) was of the interstellar medium. The astronomer Professor Hartmann determined from observations of the binary star Mintaka in Orion, that there was the element calcium in the intervening space.
Triplets
Planet Pluto was discovered by looking at photographs (i.e. 'plates' in astronomy vernacular) in a blink comparator taken with a refracting telescope, an astrograph with a 3 element 13-inch lens.
List of the largest refracting telescopes
Examples of some of the largest achromatic refracting telescopes, over diameter.
Great Paris Exhibition Telescope of 1900 () – dismantled after exhibition
Yerkes Observatory ()
Swedish 1-m Solar Telescope ()
Lick Observatory ()
Paris Observatory Meudon Great Refractor (, + )
Potsdam Great Refractor (, + )
Nice Observatory ()
John Wall () dialyte refracting telescope - the largest refractor built by an individual, at Hanwell Community Observatory
28-inch Grubb Refractor at Royal Greenwich Observatory, () aperture lens
Great Refractor of Vienna Observatory, ()
Archenhold Observatory – the longest refracting telescope ever built ( × focal length)
United States Naval Observatory refractor, ()
Newall refractor at the National Observatory of Athens ()
Lowell Observatory ()
| Technology | Telescope | null |
266611 | https://en.wikipedia.org/wiki/Optical%20telescope | Optical telescope | An optical telescope is a telescope that gathers and focuses light mainly from the visible part of the electromagnetic spectrum, to create a magnified image for direct visual inspection, to make a photograph, or to collect data through electronic image sensors.
There are three primary types of optical telescope:
Refracting telescopes, which use lenses and less commonly also prisms (dioptrics)
Reflecting telescopes, which use mirrors (catoptrics)
Catadioptric telescopes, which combine lenses and mirrors
An optical telescope's ability to resolve small details is directly related to the diameter (or aperture) of its objective (the primary lens or mirror that collects and focuses the light), and its light-gathering power is related to the area of the objective. The larger the objective, the more light the telescope collects and the finer detail it resolves.
People use optical telescopes (including monoculars and binoculars) for outdoor activities such as observational astronomy, ornithology, pilotage, hunting and reconnaissance, as well as indoor/semi-outdoor activities such as watching performance arts and spectator sports.
History
The telescope is more a discovery of optical craftsmen than an invention of a scientist. The lens and the properties of refracting and reflecting light had been known since antiquity, and theory on how they worked was developed by ancient Greek philosophers, preserved and expanded on in the medieval Islamic world, and had reached a significantly advanced state by the time of the telescope's invention in early modern Europe. But the most significant step cited in the invention of the telescope was the development of lens manufacture for spectacles, first in Venice and Florence in the thirteenth century, and later in the spectacle making centers in both the Netherlands and Germany. It is in the Netherlands in 1608 where the first documents describing a refracting optical telescope surfaced in the form of a patent filed by spectacle maker Hans Lippershey, followed a few weeks later by claims by Jacob Metius, and a third unknown applicant, that they also knew of this "art".
Word of the invention spread fast and Galileo Galilei, on hearing of the device, was making his own improved designs within a year and was the first to publish astronomical results using a telescope. Galileo's telescope used a convex objective lens and a concave eye lens, a design is now called a Galilean telescope. Johannes Kepler proposed an improvement on the design that used a convex eyepiece, often called the Keplerian Telescope.
The next big step in the development of refractors was the advent of the Achromatic lens in the early 18th century, which corrected the chromatic aberration in Keplerian telescopes up to that time—allowing for much shorter instruments with much larger objectives.
For reflecting telescopes, which use a curved mirror in place of the objective lens, theory preceded practice. The theoretical basis for curved mirrors behaving similar to lenses was probably established by Alhazen, whose theories had been widely disseminated in Latin translations of his work. Soon after the invention of the refracting telescope, Galileo, Giovanni Francesco Sagredo, and others, spurred on by their knowledge that curved mirrors had similar properties to lenses, discussed the idea of building a telescope using a mirror as the image forming objective. The potential advantages of using parabolic mirrors (primarily a reduction of spherical aberration with elimination of chromatic aberration) led to several proposed designs for reflecting telescopes, the most notable of which was published in 1663 by James Gregory and came to be called the Gregorian telescope, but no working models were built. Isaac Newton has been generally credited with constructing the first practical reflecting telescopes, the Newtonian telescope, in 1668 although due to their difficulty of construction and the poor performance of the speculum metal mirrors used it took over 100 years for reflectors to become popular. Many of the advances in reflecting telescopes included the perfection of parabolic mirror fabrication in the 18th century, silver coated glass mirrors in the 19th century, long-lasting aluminum coatings in the 20th century, segmented mirrors to allow larger diameters, and active optics to compensate for gravitational deformation. A mid-20th century innovation was catadioptric telescopes such as the Schmidt camera, which uses both a lens (corrector plate) and mirror as primary optical elements, mainly used for wide field imaging without spherical aberration.
The late 20th century has seen the development of adaptive optics and space telescopes to overcome the problems of astronomical seeing.
The electronics revolution of the early 21st century led to the development of computer-connected telescopes in the 2010s that allow non-professional skywatchers to observe stars and satellites using relatively low-cost equipment by taking advantage of digital astrophotographic techniques developed by professional astronomers over previous decades. An electronic connection to a computer (smartphone, pad, or laptop) is required to make astronomical observations from the telescopes. The digital technology allows multiple images to be stacked while subtracting the noise component of the observation producing images of Messier objects and faint stars as dim as an apparent magnitude of 15 with consumer-grade equipment.
Principles
The basic scheme is that the primary light-gathering element, the objective (1) (the convex lens or concave mirror used to gather the incoming light), focuses that light from the distant object (4) to a focal plane where it forms a real image (5). This image may be recorded or viewed through an eyepiece (2), which acts like a magnifying glass. The eye (3) then sees an inverted, magnified virtual image (6) of the object.
Inverted images
Most telescope designs produce an inverted image at the focal plane; these are referred to as inverting telescopes. In fact, the image is both turned upside down and reversed left to right, so that altogether it is rotated by 180 degrees from the object orientation. In astronomical telescopes the rotated view is normally not corrected, since it does not affect how the telescope is used. However, a mirror diagonal is often used to place the eyepiece in a more convenient viewing location, and in that case the image is erect, but still reversed left to right. In terrestrial telescopes such as spotting scopes, monoculars and binoculars, prisms (e.g., Porro prisms) or a relay lens between objective and eyepiece are used to correct the image orientation. There are telescope designs that do not present an inverted image such as the Galilean refractor and the Gregorian reflector. These are referred to as erecting telescopes.
Design variants
Many types of telescope fold or divert the optical path with secondary or tertiary mirrors. These may be integral part of the optical design (Newtonian telescope, Cassegrain reflector or similar types), or may simply be used to place the eyepiece or detector at a more convenient position. Telescope designs may also use specially designed additional lenses or mirrors to improve image quality over a larger field of view.
Characteristics
Design specifications relate to the characteristics of the telescope and how it performs optically. Several properties of the specifications may change with the equipment or accessories used with the telescope; such as Barlow lenses, star diagonals and eyepieces. These interchangeable accessories do not alter the specifications of the telescope, however they alter the way the telescope's properties function, typically magnification, apparent field of view (FOV) and actual field of view.
Surface resolvability
The smallest resolvable surface area of an object, as seen through an optical telescope, is the limited physical area that can be resolved. It is analogous to angular resolution, but differs in definition: instead of separation ability between point-light sources it refers to the physical area that can be resolved. A familiar way to express the characteristic is the resolvable ability of features such as Moon craters or Sun spots. Expression using the formula is given by twice the resolving power over aperture diameter multiplied by the objects diameter multiplied by the constant all divided by the objects apparent diameter .
Resolving power is derived from the wavelength using the same unit as aperture; where 550 nm to mm is given by: .
The constant is derived from radians to the same unit as the object's apparent diameter; where the Moon's apparent diameter of radians to arcsecs is given by: .
An example using a telescope with an aperture of 130 mm observing the Moon in a 550 nm wavelength, is given by:
The unit used in the object diameter results in the smallest resolvable features at that unit. In the above example they are approximated in kilometers resulting in the smallest resolvable Moon craters being 3.22 km in diameter. The Hubble Space Telescope has a primary mirror aperture of 2400 mm that provides a surface resolvability of Moon craters being 174.9 meters in diameter, or sunspots of 7365.2 km in diameter.
Angular resolution
Ignoring blurring of the image by turbulence in the atmosphere (atmospheric seeing) and optical imperfections of the telescope, the angular resolution of an optical telescope is determined by the diameter of the primary mirror or lens gathering the light (also termed its "aperture").
The Rayleigh criterion for the resolution limit (in radians) is given by
where is the wavelength and is the aperture. For visible light ( = 550 nm) in the small-angle approximation, this equation can be rewritten:
Here, denotes the resolution limit in arcseconds and is in millimeters.
In the ideal case, the two components of a double star system can be discerned even if separated by slightly less than . This is taken into account by the Dawes limit
The equation shows that, all else being equal, the larger the aperture, the better the angular resolution. The resolution is not given by the maximum magnification (or "power") of a telescope. Telescopes marketed by giving high values of the maximum power often deliver poor images.
For large ground-based telescopes, the resolution is limited by atmospheric seeing. This limit can be overcome by placing the telescopes above the atmosphere, e.g., on the summits of high mountains, on balloons and high-flying airplanes, or in space. Resolution limits can also be overcome by adaptive optics, speckle imaging or lucky imaging for ground-based telescopes.
Recently, it has become practical to perform aperture synthesis with arrays of optical telescopes. Very high resolution images can be obtained with groups of widely spaced smaller telescopes, linked together by carefully controlled optical paths, but these interferometers can only be used for imaging bright objects such as stars or measuring the bright cores of active galaxies.
Focal length and focal ratio
The focal length of an optical system is a measure of how strongly the system converges or diverges light. For an optical system in air, it is the distance over which initially collimated rays are brought to a focus. A system with a shorter focal length has greater optical power than one with a long focal length; that is, it bends the rays more strongly, bringing them to a focus in a shorter distance. In astronomy, the f-number is commonly referred to as the focal ratio notated as . The focal ratio of a telescope is defined as the focal length of an objective divided by its diameter or by the diameter of an aperture stop in the system. The focal length controls the field of view of the instrument and the scale of the image that is presented at the focal plane to an eyepiece, film plate, or CCD.
An example of a telescope with a focal length of 1200 mm and aperture diameter of 254 mm is given by:
Numerically large Focal ratios are said to be long or slow. Small numbers are short or fast. There are no sharp lines for determining when to use these terms, and an individual may consider their own standards of determination. Among contemporary astronomical telescopes, any telescope with a focal ratio slower (bigger number) than f/12 is generally considered slow, and any telescope with a focal ratio faster (smaller number) than f/6, is considered fast. Faster systems often have more optical aberrations away from the center of the field of view and are generally more demanding of eyepiece designs than slower ones. A fast system is often desired for practical purposes in astrophotography with the purpose of gathering more photons in a given time period than a slower system, allowing time lapsed photography to process the result faster.
Wide-field telescopes (such as astrographs), are used to track satellites and asteroids, for cosmic-ray research, and for astronomical surveys of the sky. It is more difficult to reduce optical aberrations in telescopes with low f-ratio than in telescopes with larger f-ratio.
Light-gathering power
The light-gathering power of an optical telescope, also referred to as light grasp or aperture gain, is the ability of a telescope to collect a lot more light than the human eye. Its light-gathering power is probably its most important feature. The telescope acts as a light bucket, collecting all of the photons that come down on it from a far away object, where a larger bucket catches more photons resulting in more received light in a given time period, effectively brightening the image. This is why the pupils of your eyes enlarge at night so that more light reaches the retinas. The gathering power compared against a human eye is the squared result of the division of the aperture over the observer's pupil diameter , with an average adult having a pupil diameter of 7 mm. Younger persons host larger diameters, typically said to be 9 mm, as the diameter of the pupil decreases with age.
An example gathering power of an aperture with 254 mm compared to an adult pupil diameter being 7 mm is given by:
Light-gathering power can be compared between telescopes by comparing the areas of the two different apertures.
As an example, the light-gathering power of a 10-meter telescope is 25x that of a 2-meter telescope:
For a survey of a given area, the field of view is just as important as raw light gathering power. Survey telescopes such as the Large Synoptic Survey Telescope try to maximize the product of mirror area and field of view (or etendue) rather than raw light gathering ability alone.
Magnification
The magnification through a telescope makes an object appear larger while limiting the FOV. Magnification is often misleading as the optical power of the telescope, its characteristic is the most misunderstood term used to describe the observable world. At higher magnifications the image quality significantly reduces, usage of a Barlow lens increases the effective focal length of an optical system—multiplies image quality reduction.
Similar minor effects may be present when using star diagonals, as light travels through a multitude of lenses that increase or decrease effective focal length. The quality of the image generally depends on the quality of the optics (lenses) and viewing conditions—not on magnification.
Magnification itself is limited by optical characteristics. With any telescope or microscope, beyond a practical maximum magnification, the image looks bigger but shows no more detail. It occurs when the finest detail the instrument can resolve is magnified to match the finest detail the eye can see. Magnification beyond this maximum is sometimes called empty magnification.
To get the most detail out of a telescope, it is critical to choose the right magnification for the object being observed. Some objects appear best at low power, some at high power, and many at a moderate magnification. There are two values for magnification, a minimum and maximum. A wider field of view eyepiece may be used to keep the same eyepiece focal length whilst providing the same magnification through the telescope. For a good quality telescope operating in good atmospheric conditions, the maximum usable magnification is limited by diffraction.
Visual
The visual magnification of the field of view through a telescope can be determined by the telescope's focal length divided by the eyepiece focal length (or diameter). The maximum is limited by the focal length of the eyepiece.
An example of visual magnification using a telescope with a 1200 mm focal length and 3 mm eyepiece is given by:
Minimum
There are two issues constraining the lowest useful magnification on a telescope:
The light beam exiting the eyepiece needs to be small enough to enter the pupil of the observer's eye. If the cylinder of light emerging from they eyepiece is too wide to enter the observer's eye, some of the light gathered by the telescope will be wasted, and the image seen will be dimmer and less clear than it would be at a higher magnification.
For telescope designs with obstructions in the light path (e.g. most catadioptric telescopes, but not spyglass-style refracting telescopes) the magnification must be high enough to keep the central obstruction out of focus, to prevent it from coming into view as a central "black spot". Both of these issues depend on the size of the pupil of the observer's eye, which will be narrower in daylight and wider in the dark.
Both constraints boil down to approximately the same rule: The magnification of the viewed image, must be high enough to make the eyepiece exit pupil, no larger than the pupil of the observer's own eye. The formula for the eypiece exit pupil is
where is the light-collecting diameter of the telescope's aperture.
Dark-adapted pupil sizes range from 8–9 mm for young children, to a "normal" or standard value of 7 mm for most adults aged 30–40, to 5–6 mm for retirees in their 60s and 70s. A lifetime spent exposed to chronically bright ambient light, such as sunlight reflected off of open fields of snow, or white-sand beaches, or cement, will tend to make individuals' pupils permanently smaller. Sunglasses greatly help, but once shrunk by long-time over-exposure to bright light, even the use of opthamalogic drugs cannot restore lost pupil size. Most observers' eyes instantly respond to darkness by widening the pupil to almost its maximum, although complete adaption to night vision generally takes at least a half-hour. (There is usually a slight extra widening of the pupil the longer the pupil remains dilated / relaxed.)
The improvement in brightness with reduced magnification has a limit related to something called the exit pupil. The exit pupil is the cylinder of light exiting the eyepiece and entering the pupil of the eye; hence the lower the magnification, the larger the exit pupil. It is the image of the shrunken sky-viewing aperture of the telescope, reduced by the magnification factor, of the eyepiece-telescope combination:
where is the focal length of the telescope and is the focal length of the eyepiece.
Ideally, the exit pupil of the eyepiece, matches the pupil of the observer's eye: If the exit pupil from the eyepiece is larger than the pupil of individual observer's eye, some of the light delivered from the telescope will be cut off. If the eyepiece exit pupil is the same or smaller than the pupil of the observer's eye, then all of the light collected by the telescope aperture will enter the eye, with lower magnification producing a brighter image, as long as all of the captured light gets into the eye.
The minimum can be calculated by dividing the telescope aperture over the largest tolerated exit pupil diameter
Decreasing the magnification past this limit will not increase brightness nor improve clarity: Beyond this limit there is no benefit from lower magnification. Likewise calculating the exit pupil is a division of the aperture diameter and the visual magnification used. The minimum often may not be reachable with some telescopes, a telescope with a very long focal length may require a longer focal length eyepiece than is available.
An example of the lowest usable magnification using a fairly common 10″ (254 mm) aperture and the standard adult 7 mm maximum exit pupil is given by: If the telescope happened to have a focal length (), the longest recommended eyepiece focal length () would be An eyepiece of the same apparent field-of-view but longer focal-length will deliver a wider true field of view, but dimmer image. If the telescope has a central obstruction (e.g. a Newtonian, Maksutov, or Schmidt–Cassegrain telescope) it is also likely that the low magnification will make the obstruction come into focus enough to make a black spot in the middle of the image.
Calculating in the other direction, the exit pupil diameter of a 254 mm telescope aperture at 60× magnification is given by: well within pupil size of dark-adapted eyes of observers of almost all ages. Assuming the same telescope focal length as above, the eyepiece focal length that would produce a 60× magnification is
Optimum
The following are rules-of-thumb for useful magnifications appropriate to different type objects:
For small objects with low surface brightness (such as galaxies), use a moderate magnification.
For small objects with moderate surface brightness (such as planetary nebulae), use a high magnification.
For small objects with high surface brightness (such as planets), use the highest magnification that the current night's "seeing" will allow, and consider adding in astronomical filters to sharpen the image.
For large objects (such as the Andromeda Galaxy or wide-field diffuse nebulae), regardless of surface brightness use low magnification, often in the range of minimum magnification.
For very to extremely bright, large objects (the Moon and the Sun) narrow-down the aperture of the telescope by covering it with a piece of cardboard with a small hole in it, and insert filters as-needed to both cut down excess brightness and to enhance the contrast of surface features.
Only personal experience determines the best optimum magnifications for objects, relying on observational skills and seeing conditions, and the status of the pupil of observer's eye at the moment (e.g. a lower magnification may be required if there is enough moonlight to prevent complete dark adaption).
Field of view
Field of view is the extent of the observable world seen at any given moment, through an instrument (e.g., telescope or binoculars), or by naked eye. There are various expressions of field of view, being a specification of an eyepiece or a characteristic determined from an eyepiece and telescope combination. A physical limit derives from the combination where the FOV cannot be viewed larger than a defined maximum, due to diffraction of the optics.
Apparent
Apparent field of view (commonly referred to as AFOV) is the perceived angular size of the field stop of the eyepiece, typically measured in degrees. It is a fixed property of the eyepiece's optical design, with common commercially available eyepieces offering a range of apparent fields from 40° to 120°. The apparent field of view of an eyepiece is limited by a combination of the eyepiece's field stop diameter, and focal length, and is independent of magnification used.
In an eyepiece with a very wide apparent field of view, the observer may perceive that the view through the telescope stretches out to their peripheral vision, giving a sensation that they are no longer looking through an eyepiece, or that they are closer to the subject of interest than they really are. In contrast, an eyepiece with a narrow apparent field of view may give the sensation of looking through a tunnel or small porthole window, with the black field stop of the eyepiece occupying most of the observer's vision.
A wider apparent field of view permits the observer to see more of the subject of interest (that is, a wider true field of view) without reducing magnification to do so. However, the relationship between true field of view, apparent field of view, and magnification is not direct, due to increasing distortion characteristics that correlate with wider apparent fields of view. Instead, both true field of view and apparent field of view are consequences of the eyepiece's field stop diameter.
Apparent field of view differs from true field of view in so far as true field of view varies with magnification, whereas apparent field of view does not. The wider field stop of a wide angle eyepiece permits the viewing of a wider section of the real image formed at the telescope's focal plane, thus impacting the calculated true field of view.
An eyepiece's apparent field of view can influence total view brightness as perceived by the eye, since the apparent angular size of the field stop will determine how much of the observer's retina is illuminated by the exit pupil formed by the eyepiece. However, apparent field of view has no impact on the apparent surface brightness (that is, brightness per unit area) of objects contained within the field of view.
True
True FOV is the width of what is actually seen through any given eyepiece / telescope combination.
There are two formulae for calculating true field of view:
Apparent field of view method given by , where is the true FOV, is the apparent field of view of the eyepiece, and is the magnification being used.
Eyepiece field stop method given by , where is the true FOV, is the eyepiece field stop diameter in millimeters and is the focal length of the telescope in millimeters.
The eyepiece field stop method is more accurate than the apparent field of view method, however not all eyepieces have an easily knowable field stop diameter.
Maximum
Max FOV is the maximum useful true field of view limited by the optics of the telescope. It is a physical limitation where increases beyond the maximum remain at maximum. Max FOV is the barrel size over the telescope's focal length converted from radian to degrees.
An example of max FOV using a telescope with a barrel size of 31.75 mm (1.25 inches) and focal length of 1200 mm is given by:
Observing through a telescope
There are many properties of optical telescopes and the complexity of observation using one can be a daunting task; experience and experimentation are the major contributors to understanding how to maximize one's observations. In practice, only two main properties of a telescope determine how observation differs: the focal length and aperture. These relate as to how the optical system views an object or range and how much light is gathered through an ocular eyepiece. Eyepieces further determine how the field of view and magnification of the observable world change.
Observable world
The observable world is what can be seen using a telescope. When viewing an object or range, the observer may use many different techniques. Understanding what can be viewed and how to view it depends on the field of view. Viewing an object at a size that fits entirely in the field of view is measured using the two telescope properties—focal length and aperture, with the inclusion of an ocular eyepiece with suitable focal length (or diameter). Comparing the observable world and the angular diameter of an object shows how much of the object we see. However, the relationship with the optical system may not result in high surface brightness. Celestial objects are often dim because of their vast distance, and detail may be limited by diffraction or unsuitable optical properties.
Field of view and magnification relationship
Finding what can be seen through the optical system begins with the eyepiece providing the field of view and magnification; the magnification is given by the division of the telescope and eyepiece focal lengths. Using an example of an amateur telescope such as a Newtonian telescope with an aperture of 130 mm (5") and focal length of 650 mm (25.5 inches), one uses an eyepiece with a focal length of 8 mm and apparent FOV of 52°. The magnification at which the observable world is viewed is given by: . The field of view requires the magnification, which is formulated by its division over the apparent field of view: . The resulting true field of view is 0.64°, not allowing an object such as the Orion nebula, which appears elliptical with an angular diameter of 65 × 60 arcminutes, to be viewable through the telescope in its entirety, where the whole of the nebula is within the observable world. Using methods such as this can greatly increase one's viewing potential ensuring the observable world can contain the entire object, or whether to increase or decrease magnification viewing the object in a different aspect.
Brightness factor
The surface brightness at such a magnification significantly reduces, resulting in a far dimmer appearance. A dimmer appearance results in less visual detail of the object. Details such as matter, rings, spiral arms, and gases may be completely hidden from the observer, giving a far less complete view of the object or range. Physics dictates that at the theoretical minimum magnification of the telescope, the surface brightness is at 100%. Practically, however, various factors prevent 100% brightness; these include telescope limitations (focal length, eyepiece focal length, etc.) and the age of the observer.
Age plays a role in brightness, as a contributing factor is the observer's pupil. With age the pupil naturally shrinks in diameter; generally accepted a young adult may have a 7 mm diameter pupil, an older adult as little as 5 mm, and a younger person larger at 9 mm. The minimum magnification can be expressed as the division of the aperture and pupil diameter given by: . A problematic instance may be apparent, achieving a theoretical surface brightness of 100%, as the required effective focal length of the optical system may require an eyepiece with too large a diameter.
Some telescopes cannot achieve the theoretical surface brightness of 100%, while some telescopes can achieve it using a very small-diameter eyepiece. To find what eyepiece is required to get minimum magnification one can rearrange the magnification formula, where it is now the division of the telescope's focal length over the minimum magnification: . An eyepiece of 35 mm is a non-standard size and would not be purchasable; in this scenario
to achieve 100% one would require a standard manufactured eyepiece size of 40 mm. As the eyepiece has a larger focal length than the minimum magnification, an abundance of wasted light is not received through the eyes.
Exit pupil
The limit to the increase in surface brightness as one reduces magnification is the exit pupil: a cylinder of light that projects out the eyepiece to the observer. An exit pupil must match or be smaller in diameter than one's pupil to receive the full amount of projected light; a larger exit pupil results in the wasted light. The exit pupil can be derived with from division of the telescope aperture and the minimum magnification , derived by: . The pupil and exit pupil are almost identical in diameter, giving no wasted observable light with the optical system. A 7 mm pupil falls slightly short of 100% brightness, where the surface brightness can be measured from the product of the constant 2, by the square of the pupil resulting in: . The limitation here is the pupil diameter; it is an unfortunate result and degrades with age. Some observable light loss is expected and decreasing the magnification cannot increase surface brightness once the system has reached its minimum usable magnification, hence why the term is referred to as usable.
Image Scale
When using a CCD to record observations, the CCD is placed in the focal plane. Image scale (sometimes called plate scale) is how the angular size of the object being observed is related to the physical size of the projected image in the focal plane
where is the image scale, is the angular size of the observed object, and is the physical size of the projected image. In terms of focal length image scale is
where is measured in radians per meter (rad/m), and is measured in meters. Normally is given in units of arcseconds per millimeter ("/mm). So if the focal length is measured in millimeters, the image scale is
The derivation of this equation is fairly straightforward and the result is the same for reflecting or refracting telescopes. However, conceptually it is easier to derive by considering a reflecting telescope. If an extended object with angular size is observed through a telescope, then due to the Laws of reflection and Trigonometry the size of the image projected onto the focal plane will be
The image scale (angular size of object divided by size of projected image) will be
and by using the small angle relation , when (N.B. only valid if is in radians), we obtain
Imperfect images
No telescope can form a perfect image. Even if a reflecting telescope could have a perfect mirror, or a refracting telescope could have a perfect lens, the effects of aperture diffraction are unavoidable. In reality, perfect mirrors and perfect lenses do not exist, so image aberrations in addition to aperture diffraction must be taken into account. Image aberrations can be broken down into two main classes, monochromatic, and polychromatic. In 1857, Philipp Ludwig von Seidel (1821–1896) decomposed the first order monochromatic aberrations into five constituent aberrations. They are now commonly referred to as the five Seidel Aberrations.
The five Seidel aberrations
Spherical aberration The difference in focal length between paraxial rays and marginal rays, proportional to the square of the objective diameter.
Coma A defect by which points appear as comet-like asymmetrical patches of light with tails, which makes measurement very imprecise. Its magnitude is usually deduced from the optical sine theorem.
Astigmatism The image of a point forms focal lines at the sagittal and tangental foci and in between (in the absence of coma) an elliptical shape.
Petzval field curvature The Petzval field curvature means that the image, instead of lying in a plane, actually lies on a curved surface, described as hollow or round. This causes problems when a flat imaging device is used e.g., a photographic plate or CCD image sensor.
Distortion Either barrel or pincushion, a radial distortion that must be corrected when combining multiple images (similar to stitching multiple photos into a panoramic photo).
Optical defects are always listed in the above order, since this expresses their interdependence as first order aberrations via moves of the exit/entrance pupils. The first Seidel aberration, Spherical Aberration, is independent of the position of the exit pupil (as it is the same for axial and extra-axial pencils). The second, coma, changes as a function of pupil distance and spherical aberration, hence the well-known result that it is impossible to correct the coma in a lens free of spherical aberration by simply moving the pupil. Similar dependencies affect the remaining aberrations in the list.
Chromatic aberrations
Longitudinal chromatic aberration: As with spherical aberration this is the same for axial and oblique pencils.
Transverse chromatic aberration (chromatic aberration of magnification)
Astronomical research telescopes
Optical telescopes have been used in astronomical research since the time of their invention in the early 17th century. Many types have been constructed over the years depending on the optical technology, such as refracting and reflecting, the nature of the light or object being imaged, and even where they are placed, such as space telescopes. Some are classified by the task they perform such as solar telescopes.
Large reflectors
Nearly all large research-grade astronomical telescopes are reflectors. Some reasons are:
In a lens the entire volume of material has to be free of imperfection and inhomogeneities, whereas in a mirror, only one surface has to be perfectly polished.
Light of different colors travels through a medium other than vacuum at different speeds. This causes chromatic aberration.
Reflectors work in a wider spectrum of light since certain wavelengths are absorbed when passing through glass elements like those found in a refractor or catadioptric.
There are technical difficulties involved in manufacturing and manipulating large-diameter lenses. One of them is that all real materials sag in gravity. A lens can only be held by its perimeter. A mirror, on the other hand, can be supported by the whole side opposite to its reflecting face.
Most large research reflectors operate at different focal planes, depending on the type and size of the instrument being used. These including the prime focus of the main mirror, the cassegrain focus (light bounced back down behind the primary mirror), and even external to the telescope all together (such as the Nasmyth and coudé focus).
A new era of telescope making was inaugurated by the Multiple Mirror Telescope (MMT), with a mirror composed of six segments synthesizing a mirror of 4.5 meters diameter. This has now been replaced by a single 6.5 m mirror. Its example was followed by the Keck telescopes with 10 m segmented mirrors.
The largest current ground-based telescopes have a primary mirror of between 6 and 11 meters in diameter. In this generation of telescopes, the mirror is usually very thin, and is kept in an optimal shape by an array of actuators (see active optics). This technology has driven new designs for future telescopes with diameters of 30, 50 and even 100 meters.
Relatively cheap, mass-produced ~2 meter telescopes have recently been developed and have made a significant impact on astronomy research. These allow many astronomical targets to be monitored continuously, and for large areas of sky to be surveyed. Many are robotic telescopes, computer controlled over the internet (see e.g. the Liverpool Telescope and the Faulkes Telescope North and South), allowing automated follow-up of astronomical events.
Initially the detector used in telescopes was the human eye. Later, the sensitized photographic plate took its place, and the spectrograph was introduced, allowing the gathering of spectral information. After the photographic plate, successive generations of electronic detectors, such as the charge-coupled device (CCDs), have been perfected, each with more sensitivity and resolution, and often with a wider wavelength coverage.
Current research telescopes have several instruments to choose from such as:
imagers, of different spectral responses
spectrographs, useful in different regions of the spectrum
polarimeters, that detect light polarization.
The phenomenon of optical diffraction sets a limit to the resolution and image quality that a telescope can achieve, which is the effective area of the Airy disc, which limits how close two such discs can be placed. This absolute limit is called the diffraction limit (and may be approximated by the Rayleigh criterion, Dawes limit or Sparrow's resolution limit). This limit depends on the wavelength of the studied light (so that the limit for red light comes much earlier than the limit for blue light) and on the diameter of the telescope mirror. This means that a telescope with a certain mirror diameter can theoretically resolve up to a certain limit at a certain wavelength. For conventional telescopes on Earth, the diffraction limit is not relevant for telescopes bigger than about 10 cm. Instead, the seeing, or blur caused by the atmosphere, sets the resolution limit. But in space, or if adaptive optics are used, then reaching the diffraction limit is sometimes possible. At this point, if greater resolution is needed at that wavelength, a wider mirror has to be built or aperture synthesis performed using an array of nearby telescopes.
In recent years, a number of technologies to overcome the distortions caused by atmosphere on ground-based telescopes have been developed, with good results. See adaptive optics, speckle imaging and optical interferometry.
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266861 | https://en.wikipedia.org/wiki/Reflecting%20telescope | Reflecting telescope | A reflecting telescope (also called a reflector) is a telescope that uses a single or a combination of curved mirrors that reflect light and form an image. The reflecting telescope was invented in the 17th century by Isaac Newton as an alternative to the refracting telescope which, at that time, was a design that suffered from severe chromatic aberration. Although reflecting telescopes produce other types of optical aberrations, it is a design that allows for very large diameter objectives. Almost all of the major telescopes used in astronomy research are reflectors. Many variant forms are in use and some employ extra optical elements to improve image quality or place the image in a mechanically advantageous position. Since reflecting telescopes use mirrors, the design is sometimes referred to as a catoptric telescope.
From the time of Newton to the 1800s, the mirror itself was made of metalusually speculum metal. This type included Newton's first designs and the largest telescope of the 19th century, the Leviathan of Parsonstown with a wide metal mirror. In the 19th century a new method using a block of glass coated with very thin layer of silver began to become more popular by the turn of the century. Common telescopes which led to the Crossley and Harvard reflecting telescopes, which helped establish a better reputation for reflecting telescopes as the metal mirror designs were noted for their drawbacks. Chiefly the metal mirrors only reflected about of the light and the metal would tarnish. After multiple polishings and tarnishings, the mirror could lose its precise figuring needed.
Reflecting telescopes became extraordinarily popular for astronomy and many famous telescopes, such as the Hubble Space Telescope, and popular amateur models use this design. In addition, the reflection telescope principle was applied to other electromagnetic wavelengths, and for example, X-ray telescopes also use the reflection principle to make image-forming optics.
History
The idea that curved mirrors behave like lenses dates back at least to Alhazen's 11th century treatise on optics, works that had been widely disseminated in Latin translations in early modern Europe. Soon after the invention of the refracting telescope, Galileo, Giovanni Francesco Sagredo, and others, spurred on by their knowledge of the principles of curved mirrors, discussed the idea of building a telescope using a mirror as the image forming objective. There were reports that the Bolognese Cesare Caravaggi had constructed one around 1626 and the Italian professor Niccolò Zucchi, in a later work, wrote that he had experimented with a concave bronze mirror in 1616, but said it did not produce a satisfactory image. The potential advantages of using parabolic mirrors, primarily reduction of spherical aberration with no chromatic aberration, led to many proposed designs for reflecting telescopes. These included one by James Gregory, published in 1663. In 1673, experimental scientist Robert Hooke was able to build this type of telescope, which became known as the Gregorian telescope.
Five years after Gregory designed his telescope and five years before Hooke built the first such Gregorian telescope, Isaac Newton in 1668 built his own reflecting telescope, which is generally acknowledged as the first reflecting telescope. It used a spherically ground metal primary mirror and a small diagonal mirror in an optical configuration that has come to be known as the Newtonian telescope.
Despite the theoretical advantages of the reflector design, the difficulty of construction and the poor performance of the speculum metal mirrors being used at the time meant it took over 100 years for them to become popular. Many of the advances in reflecting telescopes included the perfection of parabolic mirror fabrication in the 18th century, silver coated glass mirrors in the 19th century (built by Léon Foucault in 1858), long-lasting aluminum coatings in the 20th century, segmented mirrors to allow larger diameters, and active optics to compensate for gravitational deformation. A mid-20th century innovation was catadioptric telescopes such as the Schmidt camera, which use both a spherical mirror and a lens (called a corrector plate) as primary optical elements, mainly used for wide-field imaging without spherical aberration.
The late 20th century has seen the development of adaptive optics and lucky imaging to overcome the problems of seeing, and reflecting telescopes are ubiquitous on space telescopes and many types of spacecraft imaging devices.
Technical considerations
A curved primary mirror is the reflector telescope's basic optical element that creates an image at the focal plane. The distance from the mirror to the focal plane is called the focal length. Film or a digital sensor may be located here to record the image, or a secondary mirror may be added to modify the optical characteristics and/or redirect the light to film, digital sensors, or an eyepiece for visual observation.
The primary mirror in most modern telescopes is composed of a solid glass cylinder whose front surface has been ground to a spherical or parabolic shape. A thin layer of aluminum is vacuum deposited onto the mirror, forming a highly reflective first surface mirror.
Some telescopes use primary mirrors which are made differently. Molten glass is rotated to make its surface paraboloidal, and is kept rotating while it cools and solidifies. (See Rotating furnace.) The resulting mirror shape approximates a desired paraboloid shape that requires minimal grinding and polishing to reach the exact figure needed.
Optical errors
Reflecting telescopes, just like any other optical system, do not produce "perfect" images. The need to image objects at distances up to infinity, view them at different wavelengths of light, along with the requirement to have some way to view the image the primary mirror produces, means there is always some compromise in a reflecting telescope's optical design.
Because the primary mirror focuses light to a common point in front of its own reflecting surface almost all reflecting telescope designs have a secondary mirror, film holder, or detector near that focal point partially obstructing the light from reaching the primary mirror. Not only does this cause some reduction in the amount of light the system collects, it also causes a loss in contrast in the image due to diffraction effects of the obstruction as well as diffraction spikes caused by most secondary support structures.
The use of mirrors avoids chromatic aberration but they produce other types of aberrations. A simple spherical mirror cannot bring light from a distant object to a common focus since the reflection of light rays striking the mirror near its edge do not converge with those that reflect from nearer the center of the mirror, a defect called spherical aberration. To avoid this problem most reflecting telescopes use parabolic shaped mirrors, a shape that can focus all the light to a common focus. Parabolic mirrors work well with objects near the center of the image they produce, (light traveling parallel to the mirror's optical axis), but towards the edge of that same field of view they suffer from off axis aberrations:
Coma – an aberration where point sources (stars) at the center of the image are focused to a point but typically appears as "comet-like" radial smudges that get worse towards the edges of the image.
Field curvature – The best image plane is in general curved, which may not correspond to the detector's shape and leads to a focus error across the field. It is sometimes corrected by a field flattening lens.
Astigmatism – an azimuthal variation of focus around the aperture causing point source images off-axis to appear elliptical. Astigmatism is not usually a problem in a narrow field of view, but in a wide field image it gets rapidly worse and varies quadratically with field angle.
Distortion – Distortion does not affect image quality (sharpness) but does affect object shapes. It is sometimes corrected by image processing.
There are reflecting telescope designs that use modified mirror surfaces (such as the Ritchey–Chrétien telescope) or some form of correcting lens (such as catadioptric telescopes) that correct some of these aberrations.
Use in astronomical research
Nearly all large research-grade astronomical telescopes are reflectors. There are several reasons for this:
Reflectors work in a wider spectrum of light since certain wavelengths are absorbed when passing through glass elements like those found in a refractor or in a catadioptric telescope.
In a lens the entire volume of material has to be free of imperfection and inhomogeneities, whereas in a mirror, only one surface has to be perfectly polished.
Light of different wavelengths travels through a medium other than vacuum at different speeds. This causes chromatic aberration. Reducing this to acceptable levels usually involves a combination of two or three aperture sized lenses (see achromat and apochromat for more details). The cost of such systems therefore scales significantly with aperture size. An image obtained from a mirror does not suffer from chromatic aberration to begin with, and the cost of the mirror scales much more modestly with its size.
There are structural problems involved in manufacturing and manipulating large-aperture lenses. Since a lens can only be held in place by its edge, the center of a large lens will sag due to gravity, distorting the image it produces. The largest practical lens size in a refracting telescope is around 1 meter. In contrast, a mirror can be supported by the whole side opposite its reflecting face, allowing for reflecting telescope designs that can overcome gravitational sag. The largest reflector designs currently exceed 10 meters in diameter.
Reflecting telescope designs
Gregorian
The Gregorian telescope, described by Scottish astronomer and mathematician James Gregory in his 1663 book Optica Promota, employs a concave secondary mirror that reflects the image back through a hole in the primary mirror. This produces an upright image, useful for terrestrial observations. Some small spotting scopes are still built this way. There are several large modern telescopes that use a Gregorian configuration such as the Vatican Advanced Technology Telescope, the Magellan telescopes, the Large Binocular Telescope, and the Giant Magellan Telescope.
Newtonian
The Newtonian telescope was the first successful reflecting telescope, completed by Isaac Newton in 1668. It usually has a paraboloid primary mirror but at focal ratios of about f/10 or longer a spherical primary mirror can be sufficient for high visual resolution. A flat secondary mirror reflects the light to a focal plane at the side of the top of the telescope tube. It is one of the simplest and least expensive designs for a given size of primary, and is popular with amateur telescope makers as a home-build project.
The Cassegrain design and its variations
The Cassegrain telescope (sometimes called the "Classic Cassegrain") was first published in a 1672 design attributed to Laurent Cassegrain. It has a parabolic primary mirror, and a hyperbolic secondary mirror that reflects the light back down through a hole in the primary. The folding and diverging effect of the secondary mirror creates a telescope with a long focal length while having a short tube length.
Ritchey–Chrétien
The Ritchey–Chrétien telescope, invented by George Willis Ritchey and Henri Chrétien in the early 1910s, is a specialized Cassegrain reflector which has two hyperbolic mirrors (instead of a parabolic primary). It is free of coma and spherical aberration at a nearly flat focal plane if the primary and secondary curvature are properly figured, making it well suited for wide field and photographic observations. Almost every professional reflector telescope in the world is of the Ritchey–Chrétien design.
Three-mirror anastigmat
Including a third curved mirror allows correction of the remaining distortion, astigmatism, from the Ritchey–Chrétien design. This allows much larger fields of view.
Dall–Kirkham
The Dall–Kirkham Cassegrain telescope's design was created by Horace Dall in 1928 and took on the name in an article published in Scientific American in 1930 following discussion between amateur astronomer Allan Kirkham and Albert G. Ingalls, the magazine editor at the time. It uses a concave elliptical primary mirror and a convex spherical secondary. While this system is easier to grind than a classic Cassegrain or Ritchey–Chrétien system, it does not correct for off-axis coma. Field curvature is actually less than a classical Cassegrain. Because this is less noticeable at longer focal ratios, Dall–Kirkhams are seldom faster than f/15.
Off-axis designs
There are several designs that try to avoid obstructing the incoming light by eliminating the secondary or moving any secondary element off the primary mirror's optical axis, commonly called off-axis optical systems.
Herschelian
The Herschelian reflector is named after William Herschel, who used this design to build very large telescopes including the 40-foot telescope in 1789. In the Herschelian reflector the primary mirror is tilted so the observer's head does not block the incoming light. Although this introduces geometrical aberrations, Herschel employed this design to avoid the use of a Newtonian secondary mirror since the speculum metal mirrors of that time tarnished quickly and could only achieve 60% reflectivity.
Schiefspiegler
A variant of the Cassegrain, the Schiefspiegler telescope ("skewed" or "oblique reflector") uses tilted mirrors to avoid the secondary mirror casting a shadow on the primary. However, while eliminating diffraction patterns this leads to an increase in coma and astigmatism. These defects become manageable at large focal ratios — most Schiefspieglers use f/15 or longer, which tends to restrict useful observations to objects which fit in a moderate field of view. A 6" (150mm) f/15 telescope offers a maximum 0.75 degree field of view using 1.25" eyepieces.
A number of variations are common, with varying numbers of mirrors of different types. The Kutter (named after its inventor Anton Kutter) style uses a single concave primary, a convex secondary and a plano-convex lens between the secondary mirror and the focal plane, when needed (this is the case of the catadioptric Schiefspiegler).
One variation of a multi-schiefspiegler uses a concave primary, convex secondary and a parabolic tertiary. One of the interesting aspects of some Schiefspieglers is that one of the mirrors can be involved in the light path twice — each light path reflects along a different meridional path.
Stevick-Paul
Stevick-Paul telescopes are off-axis versions of Paul 3-mirror systems with an added flat diagonal mirror.
A convex secondary mirror is placed just to the side of the light entering the telescope, and positioned afocally so as to send parallel light on to the tertiary.
The concave tertiary mirror is positioned exactly twice as far to the side of the entering beam as was the convex secondary, and its own radius of curvature distant from the secondary. Because the tertiary mirror receives parallel light from the secondary, it forms an image at its focus.
The focal plane lies within the system of mirrors, but is accessible to the eye with the inclusion of a flat diagonal.
The Stevick-Paul configuration results in all optical aberrations totaling zero to the third-order, except for the Petzval surface which is gently curved.
Yolo
The Yolo was developed by Arthur S. Leonard in the mid-1960s. Like the Schiefspiegler, it is an unobstructed, tilted reflector telescope. The original Yolo consists of a primary and secondary concave mirror, with the same curvature, and the same tilt to the main axis. Most Yolos use toroidal reflectors. The Yolo design eliminates coma, but leaves significant astigmatism, which is reduced by deformation of the secondary mirror by some form of warping harness, or alternatively, polishing a toroidal figure into the secondary.
Like Schiefspieglers, many Yolo variations have been pursued. The needed amount of toroidal shape
can be transferred entirely or partially to the primary mirror. In large focal ratios optical assemblies, both primary and secondary mirror can be left spherical and a spectacle correcting lens is added between the secondary mirror and the focal plane (catadioptric Yolo). The addition of a convex, long focus tertiary mirror leads to Leonard's Solano configuration. The Solano telescope doesn't contain any toric surfaces.
Liquid-mirror telescopes
One design of telescope uses a rotating mirror consisting of a liquid metal in a tray that is spun at constant speed. As the tray spins, the liquid forms a paraboloidal surface of essentially unlimited size. This allows making very big telescope mirrors (over 6 metres), but they are limited to use by zenith telescopes.
Focal planes
Prime focus
In a prime focus design no secondary optics are used, the image is accessed at the focal point of the primary mirror. At the focal point is some type of structure for holding a film plate or electronic detector. In the past, in very large telescopes, an observer would sit inside the telescope in an "observing cage" to directly view the image or operate a camera. Nowadays CCD cameras allow for remote operation of the telescope from almost anywhere in the world. The space available at prime focus is severely limited by the need to avoid obstructing the incoming light.
Radio telescopes often have a prime focus design. The mirror is replaced by a metal surface for reflecting radio waves, and the observer is an antenna.
Cassegrain focus
For telescopes built to the Cassegrain design or other related designs, the image is formed behind the primary mirror, at the focal point of the secondary mirror. An observer views through the rear of the telescope, or a camera or other instrument is mounted on the rear. Cassegrain focus is commonly used for amateur telescopes or smaller research telescopes. However, for large telescopes with correspondingly large instruments, an instrument at Cassegrain focus must move with the telescope as it slews; this places additional requirements on the strength of the instrument support structure, and potentially limits the movement of the telescope in order to avoid collision with obstacles such as walls or equipment inside the observatory.
Nasmyth and coudé focus
Nasmyth
The Nasmyth design is similar to the Cassegrain except the light is not directed through a hole in the primary mirror; instead, a third mirror reflects the light to the side of the telescope to allow for the mounting of heavy instruments. This is a very common design in large research telescopes.
Coudé
Adding further optics to a Nasmyth-style telescope to deliver the light (usually through the declination axis) to a fixed focus point that does not move as the telescope is reoriented gives a coudé focus (from the French word for elbow). The coudé focus gives a narrower field of view than a Nasmyth focus and is used with very heavy instruments that do not need a wide field of view. One such application is high-resolution spectrographs that have large collimating mirrors (ideally with the same diameter as the telescope's primary mirror) and very long focal lengths. Such instruments could not withstand being moved, and adding mirrors to the light path to form a coudé train, diverting the light to a fixed position to such an instrument housed on or below the observing floor (and usually built as an unmoving integral part of the observatory building) was the only option. The 60-inch Hale telescope (1.5 m), Hooker Telescope, 200-inch Hale Telescope, Shane Telescope, and Harlan J. Smith Telescope all were built with coudé foci instrumentation. The development of echelle spectrometers allowed high-resolution spectroscopy with a much more compact instrument, one which can sometimes be successfully mounted on the Cassegrain focus. Since inexpensive and adequately stable computer-controlled alt-az telescope mounts were developed in the 1980s, the Nasmyth design has generally supplanted the coudé focus for large telescopes.
Fibre-fed spectrographs
For instruments requiring very high stability, or that are very large and cumbersome, it is desirable to mount the instrument on a rigid structure, rather than moving it with the telescope. Whilst transmission of the full field of view would require a standard coudé focus, spectroscopy typically involves the measurement of only a few discrete objects, such as stars or galaxies. It is therefore feasible to collect light from these objects with optical fibers at the telescope, placing the instrument at an arbitrary distance from the telescope. Examples of fiber-fed spectrographs include the planet-hunting spectrographs HARPS or ESPRESSO.
Additionally, the flexibility of optical fibers allow light to be collected from any focal plane; for example, the HARPS spectrograph utilises the Cassegrain focus of the ESO 3.6 m Telescope, whilst the Prime Focus Spectrograph is connected to the prime focus of the Subaru telescope.
| Technology | Telescope | null |
797487 | https://en.wikipedia.org/wiki/Pieris%20brassicae | Pieris brassicae | Pieris brassicae, the large white, also called cabbage butterfly, cabbage white, cabbage moth (erroneously), or in India the large cabbage white, is a butterfly in the family Pieridae. It is a close relative of the small white, Pieris rapae.
The large white is common throughout Europe, North Africa and Asia.
Distribution
The large white is common throughout Europe, north Africa, and Asia to the Himalayas often in agricultural areas, meadows and parkland. It has managed to establish a population in South Africa and in 1995 it was predicted to spread to Australia and New Zealand.
The large white is a strong flier and the British population has been reinforced in most years by migrations from the continent. Scattered reports of the large white from the north-eastern United States (New York, Rhode Island and Maine) over the past century are of a dubious nature and indicate either accidental transport or intentional release. Such introductions threaten to establish this agricultural pest in North America.
In 2010 the butterfly was found in Nelson, New Zealand where it is known as the great white butterfly. It is classed as an unwanted pest due to the potential effect on crops. For a limited period in October 2013 the Department of Conservation offered a monetary reward for the capture of the butterfly. After two weeks, the public had captured 134 butterflies, netting $10 for each one handed in. As a result of this and other containment measures, such as over 263,000 searches in the upper South Island and the release of predatory wasps, the large white was officially declared to be eradicated from New Zealand as of December 2014.
Eggs
The large white ova are pale yellow, turning darker yellow within twenty-four hours of being oviposited. A few hours prior to hatching, they become black, the shell more transparent, and the larvae visible within.
Larvae
Large white larvae experience four moultings and five instars. The first instar follows the hatching of the egg into large white larvae. The larvae are light yellow with distinctive brown heads and have soft bodies. The larvae appear to be very hairy. Following a moulting, the larvae enter the second instar. They have tubercles covered with black hair. In the third instar, large white larvae display more activity. This instar is when the larvae are observed to eat voraciously, and cause significant amounts of damage to their host plant. At this point, they are observed to be more yellow in colour, studded with black dots. Following the third instar, the larvae go through the fourth instar, with similar appearances as the larvae of the third instar, but with more aggrandized size and feeding behaviour. The large white larvae are observed to be cylindrical, robust, and elongated by the fifth instar, yellow in colour and with bright colouration on their abdomen and thorax. They are also observed to have a grey and black head. This instar requires maximum food quality and quantity in order to aid in full development, otherwise the larva dies before becoming an adult butterfly.
Imagines
For both males and females, the wings are white with black tips on the forewings. The female also has two black spots on each forewing. The underside of each wing is a pale greenish and serves as excellent camouflage when at rest. The black markings are generally darker in the summer brood. The large white butterfly's wingspan reaches 5 to 6.5 cm on average.
Male
The upperside of the male is creamy white. The forewing is irrorated (sprinkled) with black scales at the base and along the costa for a short distance. The apex and termen above vein 2 are more or less broadly black with the inner margin of the black area containing a regular even curve. In one or two specimens a small longitudinally narrow black spot was found in interspace 3. Hindwing: uniform, irrorated with black scales at base, a large black subcostal spot before the apex, and in a few specimens indications of black scaling on the termen anteriorly. The underside of the forewing is white, slightly irrorated with black scales at the base of cell and along costa. The apex is light ochraceous brown with a large black spot in outer half of interspace 1 and another quadrate black spot at base of interspace 3. The hindwing is light ochraceous brown, closely irrorated with minute black scales. The subcostal black spot before the apex shows through from the upperside. The antennae are black and white at apex. The head, thorax, and abdomen are black, with some white hairs, where underneath is whitish.
Female
The upperside of the female is similar to that of the male, but the irroration of black scales at the bases of the wings is more extended. The black area on apex and termen of forewing is broader, its inner margin less evenly curved. A conspicuous large, black spot also exists in the outer half of interspace 1 near the base of interspace 3. On the hindwing the subcostal black spot before the apex is much larger and more prominent. The underside is similar to that of the male but the apex of the forewing and the whole surface of the hindwing is a light ochraceous yellow, not ochraceous brown. The black discal spots on forewing are much larger. The antennae, head, thorax, and abdomen of the females are the same as for the male.
Habitat
The large white butterfly's habitat consists of large, open spaces, as well as farms and vegetable gardens, because of the availability of its food source. Some favoured locations include walls, fences, tree trunks, and often their food plant. They primarily hover around these locations, which should contain both wild and cultivated crucifer, as well as oil-seed rape, cabbages, and Brussels sprouts.
Reproduction and development
Mating system
These butterflies can be polyandrous, but it is not the predominant mating system. This means that, though some female butterflies can have more than one mate, most of the large white females only have one male mate at a time through a monogamous mating system.
Two generations of butterflies are produced each year. The first brood consists of adults with a spring hatching around April. The second brood is made up of adults that hatch around July. Sometimes, a third brood can be observed farther along in the summer if the weather is warm enough.
Life cycle
Oviposition
These female butterflies oviposit in clusters on the undersides of leaves because the larvae prefer the morphology of leaf undersides over the upper surface of leaves. To oviposit, the female butterflies use the tip of the abdomen and arrange the ova in specific batches.
The pre-oviposition period, which lasts three to eight days, provides ample time for these butterflies to mate. Females tend to use their forelegs to drum on the surfaces of their intended leaves as a test of the plant's suitability for breeding. If they find a suitable surface, female large whites oviposit two to three days following copulation. They oviposit approximately six to seven times in eight days. The females can pair up to mate again approximately five or more days after the previous mating.
Choosing locations for oviposition
Females rely on visual cues, such as the colours of plants, to decide where to lay their eggs. They favour green surfaces in particular to display oviposition behaviour. This colour preference could be due to the fact that the large white's food source also acts as a host plant for oviposition.
Most females choose nectar plants like buddleia or thistles, which are green and ideal plants for the larvae. These plants, used as oviposition sites, typically contain mustard oil glucosides, whose primary function is to help the larvae survive as their essential food source. For instance, previous studies have shown that the large white larvae do not survive if the adult butterflies oviposit on a different host plant such as broad bean (Vicia faba) because this bean does not contain the proper nutrients to aid larval development.
Hatching
The large white eggs hatch approximately one week after being laid and live as a group for some time. The hatching period constitutes around two to seven hours. Upon hatching, they cause a lot of damage to the host plant by eating away at and destroying the host plant. As expected during the colder moments of the day they may appear inactive - dormant.
Behaviour
Migration
The large whites are found throughout most of Eurasia, though there are some seasonal fluctuations present due to migration. The northern populations tend to be augmented during the summer migration season from butterflies from southern areas. The large whites fly starting early spring, and keep migrating until seasons shift to autumn and the resultant cold weather. This means the large whites typically take two to three flights per butterfly reproductive season.
Large white butterfly migration patterns are typically observed only when there is a disturbance. In general, the large white butterfly's migratory patterns are atypical; normally, butterflies fly towards the poles in the spring, and towards the more temperate Equator during the fall. However, they fly in random directions, excluding north, in the spring, and there is little return migration observed. However, it has been hard to track entire migratory paths, since these butterflies can migrate more than 800 kilometres; thus, individual butterflies may not migrate the 800 kilometres, but rather that other butterflies start their migrations from where the other butterflies ended.
Hibernation
Large white broods in the north have not been seen to overwinter, nor hibernate over the winter, successfully. However, they have been observed to hibernate in the south.
Territorial behaviour
Males do not display considerable amounts of territorial behaviour. It has been suggested that this could be a reason why there is no observed significant sexual dimorphism between the male and female large white butterflies.
Ecology
Diet and food selection
Large white butterflies have a preference for what types of food plant they usually eat. Studies have shown that the preference for certain plants is reliant upon the butterflies' previous experiences. The large white butterflies, then, are shown to rely on the species of food plants, the time of experience, and the choice-situation. Thus, the large white butterflies learn what types of foods they prefer, rather than relying on their sense organs or physiological changes. In contrast, this preference for adult food plant differs from the preference of female large whites using visual cues such as plant colour to determine the best host plants for oviposition.
Plants with mustard-oil glucosides are important for this butterfly because it dictates their eating behaviours, and resultant survival rates, as specified in the section regarding oviposition. This is so beneficial for large whites because their large consumption of plants containing mustard oils is the specific reason they are so distasteful to predators, such as birds. Thus, caterpillars are protected from attack, despite them being brightly coloured; in fact, the bright colouration is to signal to predators that they taste bad.
However, there is more benefit to this species' use of mustard oil glucosides. In addition to predator protection, these glucosides belong to a class of stimuli that produce the biting responses associated with eating. Some plants contain alkaloids and steroids; these reduce and inhibit the butterflies' responsiveness to mustard oil glucosides. Thus, this utilization of mustard oil glucosides dramatically affects the behaviour of the butterfly, and the resulting food selection for survival.
The food source of the larva of the white butterfly are cabbages, radishes, and the undersides of leaves. Adults feed on flower nectar.
Predators
Large white butterflies do not have a specific group of predators. Instead, they are preyed upon by a wide range of animals, and even the occasional plant. This butterfly's main predators include birds; however, large whites can also be preyed upon by species in orders such as Hymenoptera, Hemiptera, Coleoptera, Diptera, Arachnid; some species of mammals, one of reptiles, one species of insectivorous plant, and species in amphibian orders, as well as other miscellaneous insect species. The butterflies are typically preyed upon as eggs, larvae, and imagoes.
Aposematism
Large white butterflies emit an unpleasant smell which deters predators. In addition, large whites are an aposematic species, meaning that they display warning colours, which benefits the large whites against predation. This aposematic colouration occurs in the larval, pupal, and imago stages, where toxic mustard oil glycosides from food plants are stored in the individuals' bodies. Aposematism is not entirely related to Müllerian mimicry; however, large white larvae often benefit from multiple other aposematic larvae from other species, such as the larvae of Papilio machaon.
Relationship to people
Role as pests
The crops most susceptible to P. brassicae damage in areas in Europe are those in the genus Brassica (cabbage, mustard, and their allies), particularly Brussels sprouts, cabbage, cauliflower, kohlrabi, rape, swede, and turnip. The attacks to crops are rather localized and can lead to 100% crop loss in a certain area. In addition, because of its strong inclination to migrate, adults may infest new areas that were previously free from attack. Because many of the host plants of P. brassicae are sold for consumption, damage by these butterflies can cause a great reduction of crop value. Larvae may also bore into the vegetable heads of cabbage and cauliflower and cause damage. High populations of these larvae may also skeletonise their host plants. In present-day areas such as Great Britain, P. brassicae are now less threatening as pests because of natural and chemical control reasons. However, it is still considered a pest in other European countries, in China, India, Nepal, and Russia. In fact, it is estimated to cause over 40% yield loss annually on different crop vegetables in India and Turkey.
Subspecies
Subspecies include the following:
Pieris brassicae azorensis Rebel, 1917
Pieris brassicae brassicae (Linnaeus, 1758)
Pieris brassicae catoleuca Röber, 1896
Pieris brassicae cyniphia (Turati, 1924)
Pieris brassicae cypria Verity, 1908
Pieris brassicae italorum Stauder, 1921
Pieris brassicae nepalensis Gray, 1846
Pieris brassicae ottonis Röber, 1907
Pieris brassicae subtaeniata (Turati, 1929)
Pieris brassicae vazquezi Oberthür, 1914
Pieris brassicae verna Zeller, 1924
Pieris brassicae wollastoni (Butler, 1886) †
| Biology and health sciences | Lepidoptera | null |
798370 | https://en.wikipedia.org/wiki/Computer%20cooling | Computer cooling | Computer cooling is required to remove the waste heat produced by computer components, to keep components within permissible operating temperature limits. Components that are susceptible to temporary malfunction or permanent failure if overheated include integrated circuits such as central processing units (CPUs), chipsets, graphics cards, hard disk drives, and solid state drives.
Components are often designed to generate as little heat as possible, and computers and operating systems may be designed to reduce power consumption and consequent heating according to workload, but more heat may still be produced than can be removed without attention to cooling. Use of heatsinks cooled by airflow reduces the temperature rise produced by a given amount of heat. Attention to patterns of airflow can prevent the development of hotspots. Computer fans are widely used along with heatsink fans to reduce temperature by actively exhausting hot air. There are also other cooling techniques, such as liquid cooling. All modern day processors are designed to cut out or reduce their voltage or clock speed if the internal temperature of the processor exceeds a specified limit. This is generally known as Thermal Throttling in the case of reduction of clock speeds, or Thermal Shutdown in the case of a complete shutdown of the device or system.
Cooling may be designed to reduce the ambient temperature within the case of a computer, such as by exhausting hot air, or to cool a single component or small area (spot cooling). Components commonly individually cooled include the CPU, graphics processing unit (GPU) and the northbridge.
Generators of unwanted heat
Integrated circuits (e.g. CPU and GPU) are the main generators of heat in modern computers. Heat generation can be reduced by efficient design and selection of operating parameters such as voltage and frequency, but ultimately, acceptable performance can often only be achieved by managing significant heat generation.
In operation, the temperature of a computer's components will rise until the heat transferred to the surroundings is equal to the heat produced by the component, that is, when thermal equilibrium is reached. For reliable operation, the temperature must never exceed a specified maximum permissible value unique to each component. For semiconductors, instantaneous junction temperature, rather than component case, heatsink, or ambient temperature is critical.
Cooling can be impaired by:
Dust acting as a thermal insulator and impeding airflow, thereby reducing heatsink and fan performance.
Poor airflow including turbulence due to friction against impeding components such as ribbon cables, or incorrect orientation of fans, can reduce the amount of air flowing through a case and even create localized whirlpools of hot air in the case. In some cases of equipment with bad thermal design, cooling air can easily flow out through "cooling" holes before passing over hot components; cooling in such cases can often be improved by blocking of selected holes.
Poor heat transfer due to poor thermal contact between components to be cooled and cooling devices. This can be improved by the use of thermal compounds to even out surface imperfections, or even by lapping.
Damage prevention
Because high temperatures can significantly reduce life span or cause permanent damage to components, and the heat output of components can sometimes exceed the computer's cooling capacity, manufacturers often take additional precautions to ensure that temperatures remain within safe limits. A computer with thermal sensors integrated in the CPU, motherboard, chipset, or GPU can shut itself down when high temperatures are detected to prevent permanent damage, although this may not completely guarantee long-term safe operation. Before an overheating component reaches this point, it may be "throttled" until temperatures fall below a safe point using dynamic frequency scaling technology. Throttling reduces the operating frequency and voltage of an integrated circuit or disables non-essential features of the chip to reduce heat output, often at the cost of slightly or significantly reduced performance. For desktop and notebook computers, throttling is often controlled at the BIOS level. Throttling is also commonly used to manage temperatures in smartphones and tablets, where components are packed tightly together with little to no active cooling, and with additional heat transferred from the hand of the user.
The user can also perform several tasks in order to preemptively prevent damage from happening. They can perform a visual inspection of the cooler and case fans. If any of them are not spinning correctly, it is likely that they will need to be replaced. The user should also clean the fans thoroughly, since dust and debris can increase the ambient case temperature and impact fan performance. The best way to do so is with compressed air in an open space. Another preemptive technique to prevent damage is to replace the thermal paste regularly.
Mainframes and supercomputers
As electronic computers became larger and more complex, cooling of the active components became a critical factor for reliable operation. Early vacuum-tube computers, with relatively large cabinets, could rely on natural or forced air circulation for cooling. However, solid-state devices were packed much more densely and had lower allowable operating temperatures.
Starting in 1965, IBM and other manufacturers of mainframe computers sponsored intensive research into the physics of cooling densely packed integrated circuits. Many air and liquid cooling systems were devised and investigated, using methods such as natural and forced convection, direct air impingement, direct liquid immersion and forced convection, pool boiling, falling films, flow boiling, and liquid jet impingement. Mathematical analysis was used to predict temperature rises of components for each possible cooling system geometry.
IBM developed three generations of the Thermal Conduction Module (TCM) which used a water-cooled cold plate in direct thermal contact with integrated circuit packages. Each package had a thermally conductive pin pressed onto it, and helium gas surrounded chips and heat-conducting pins. The design could remove up to 27 watts from a chip and up to 2000 watts per module, while maintaining chip package temperatures of around . Systems using TCMs were the 3081 family (1980), ES/3090 (1984) and some models of the ES/9000 (1990). In the IBM 3081 processor, TCMs allowed up to 2700 watts on a single printed circuit board while maintaining chip temperature at . Thermal conduction modules using water cooling were also used in mainframe systems manufactured by other companies including Mitsubishi and Fujitsu.
The Cray-1 supercomputer designed in 1976 had a distinctive cooling system. The machine was only in height and in diameter, and consumed up to 115 kilowatts; this is comparable to the average power consumption of a few dozen Western homes or a medium-sized car. The integrated circuits used in the machine were the fastest available at the time, using emitter-coupled logic; however, the speed was accompanied by high power consumption compared to later CMOS devices.
Heat removal was critical. Refrigerant was circulated through piping embedded in vertical cooling bars in twelve columnar sections of the machine. Each of the 1662 printed circuit modules of the machine had a copper core and was clamped to the cooling bar. The system was designed to maintain the cases of integrated circuits at no more than , with refrigerant circulating at . Final heat rejection was through a water-cooled condenser. Piping, heat exchangers, and pumps for the cooling system were arranged in an upholstered bench seat around the outside of the base of the computer. About 20 percent of the machine's weight in operation was refrigerant.
In the later Cray-2, with its more densely packed modules, Seymour Cray had trouble effectively cooling the machine using the metal conduction technique with mechanical refrigeration, so he switched to 'liquid immersion' cooling. This method involved filling the chassis of the Cray-2 with a liquid called Fluorinert. Fluorinert, as its name implies, is an inert liquid that does not interfere with the operation of electronic components. As the components came to operating temperature, the heat would dissipate into the Fluorinert, which was pumped out of the machine to a chilled water heat exchanger.
Performance per watt of modern systems has greatly improved; many more computations can be carried out with a given power consumption than was possible with the integrated circuits of the 1980s and 1990s. Recent supercomputer projects such as Blue Gene rely on air cooling, which reduces cost, complexity, and size of systems compared to liquid cooling.
Air cooling
Fans
Fans are used when natural convection is insufficient to remove heat. Fans may be fitted to the computer case or attached to CPUs, GPUs, chipsets, power supply units (PSUs), hard drives, or as cards plugged into an expansion slot. Common fan sizes include 40, 60, 80, 92, 120, and 140 mm. 200, 230, 250 and 300 mm fans are sometimes used in high-performance personal computers.
Performance of fans in chassis
A computer has a certain resistance to air flowing through the chassis and components. This is the sum of all the smaller impediments to air flow, such as the inlet and outlet openings, air filters, internal chassis, and electronic components. Fans are simple air pumps that provide pressure to the air of the inlet side relative to the output side. That pressure difference moves air through the chassis, with air flowing to areas of lower pressure.
Fans generally have two published specifications: free air flow and maximum differential pressure. Free air flow is the amount of air a fan will move with zero back-pressure. Maximum differential pressure is the amount of pressure a fan can generate when completely blocked. In between these two extremes are a series of corresponding measurements of flow versus pressure which is usually presented as a graph. Each fan model will have a unique curve, like the dashed curves in the adjacent illustration.
Parallel vis-à-vis series installation
Fans can be installed parallel to each other, in series, or a combination of both. Parallel installation would be fans mounted side by side. Series installation would be a second fan in line with another fan such as an inlet fan and an exhaust fan. To simplify the discussion, it is assumed the fans are the same model.
Parallel fans will provide double the free air flow but no additional driving pressure. Series installation, on the other hand, will double the available static pressure but not increase the free air flow rate. The adjacent illustration shows a single fan versus two fans in parallel with a maximum pressure of of water and a doubled flow rate of about .
Note that air flow changes as the square root of the pressure. Thus, doubling the pressure will only increase the flow 1.41 () times, not twice as might be assumed. Another way of looking at this is that the pressure must go up by a factor of four to double the flow rate.
To determine flow rate through a chassis, the chassis impedance curve can be measured by imposing an arbitrary pressure at the inlet to the chassis and measuring the flow through the chassis. This requires fairly sophisticated equipment. With the chassis impedance curve (represented by the solid red and black lines on the adjacent curve) determined, the actual flow through the chassis as generated by a particular fan configuration is graphically shown where the chassis impedance curve crosses the fan curve. The slope of the chassis impedance curve is a square root function, where doubling the flow rate required four times the differential pressure.
In this particular example, adding a second fan provided marginal improvement with the flow for both configurations being approximately . While not shown on the plot, a second fan in series would provide slightly better performance than the parallel installation.
Temperature vis-à-vis flow rate
The equation for required airflow through a chassis is
where
= Cubic Feet per Minute ()
= Heat Transferred (kW)
= Specific Heat of Air
= Density
= Change in Temperature (in °F)
A simple conservative rule of thumb for cooling flow requirements, discounting such effects as heat loss through the chassis walls and laminar versus turbulent flow, and accounting for the constants for specific heat and density at sea level is:
For example, a typical chassis with 500 watts of load, maximum internal temperature in a environment, i.e. a difference of :
This would be actual flow through the chassis and not the free air rating of the fan. It should also be noted that "Q", the heat transferred, is a function of the heat transfer efficiency of a CPU or GPU cooler to the airflow.
Piezoelectric pump
A "dual piezo cooling jet", patented by GE, uses vibrations to pump air through the device. The initial device is three millimetres thick and consists of two nickel discs that are connected on either side to a sliver of piezoelectric ceramics. An alternating current passed through the ceramic component causes it to expand and contract at up to 150 times per second so that the nickel discs act like a bellows. Contracted, the edges of the discs are pushed apart and suck in hot air. Expanding brings the nickel discs together, expelling the air at high velocity.
The device has no bearings and does not require a motor. It is thinner and consumes less energy than typical fans. The jet can move the same amount of air as a cooling fan twice its size while consuming half as much electricity and at lower cost.
Passive cooling
Passive heatsink cooling involves attaching a block of machined or extruded metal to the part that needs cooling. A thermal adhesive may be used. More commonly for a personal computer CPU, a clamp holds the heatsink directly over the chip, with a thermal grease or thermal pad spread between. This block has fins and ridges to increase its surface area. The heat conductivity of metal is much better than that of air, and it radiates heat better than the component that it is protecting (usually an integrated circuit or CPU).
Dust buildup between the metal fins of a heatsink gradually reduces efficiency, but can be countered with a gas duster by blowing away the dust along with any other unwanted excess material.
Passive heatsinks are commonly found on older CPUs, parts that do not get very hot (such as the chipset), low-power computers, and embedded devices. Many smartphones are used passive cooling without heatsinks.
Usually a heatsink is attached to the integrated heat spreader (IHS), essentially a large, flat plate attached to the CPU, with conduction paste layered between. This dissipates or spreads the heat locally. Unlike a heatsink, a spreader is meant to redistribute heat, not to remove it. In addition, the IHS protects the fragile CPU.
Passive cooling involves no fan noise, as convection forces move air over the heatsink.
Other techniques
Liquid immersion cooling
Another growing trend due to the increasing heat density of computers, GPUs, FPGAs, and ASICs is to immerse the entire computer or select components in a thermally, but not electrically, conductive liquid. Although rarely used for the cooling of personal computers, liquid immersion is a routine method of cooling large power distribution components such as transformers. It is also becoming popular with data centers. Personal computers cooled in this manner may not require either fans or pumps, and may be cooled exclusively by passive heat exchange between the computer hardware and the enclosure it is placed in. A heat exchanger (i.e. heater core or radiator) might still be needed though, and the piping also needs to be placed correctly.
The coolant used must have sufficiently low electrical conductivity not to interfere with the normal operation of the computer. If the liquid is somewhat electrically conductive, it may cause electrical shorts between components or traces and permanently damage them. For these reasons, it is preferred that the liquid be an insulator (dielectric) and not conduct electricity.
A wide variety of liquids exist for this purpose, including transformer oils, synthetic single-phase and dual phase dielectric coolants such as 3M Fluorinert or 3M Novec. Non-purpose oils, including cooking, motor and silicone oils, have been successfully used for cooling personal computers.
Some fluids used in immersion cooling, especially hydrocarbon based materials such as mineral oils, cooking oils, and organic esters, may degrade some common materials used in computers such as rubbers, polyvinyl chloride (PVC), and thermal greases. Therefore it is critical to review the material compatibility of such fluids prior to use. Mineral oil in particular has been found to have negative effects on PVC and rubber-based wire insulation. Thermal pastes used to transfer heat to heatsinks from processors and graphic cards has been reported to dissolve in some liquids, however with negligible impact to cooling, unless the components were removed and operated in air.
Evaporation, especially for 2-phase coolants, can pose a problem, and the liquid may require either to be regularly refilled or sealed inside the computer's enclosure. Immersion cooling can allow for extremely low PUE values of 1.05, vs air cooling's 1.35, and allow for up to 100 KW of computing power (heat dissipation, TDP) per 19-inch rack, as opposed to air cooling, which usually handles up to 23 KW.
Waste heat reduction
Where powerful computers with many features are not required, less powerful computers or ones with fewer features can be used. a VIA EPIA motherboard with CPU typically dissipates approximately 25 watts of heat, whereas a more capable Pentium 4 motherboard and CPU typically dissipates around 140 watts. Computers can be powered with direct current from an external power supply unit which does not generate heat inside the computer case. The replacement of cathode-ray-tube (CRT) displays by more efficient thin-screen liquid crystal display (LCD) ones in the early twenty-first century has reduced power consumption significantly.
Heat-sinks
A component may be fitted in good thermal contact with a heatsink, a passive device with large thermal capacity and with a large surface area relative to its volume. Heatsinks are usually made of a metal with high thermal conductivity such as aluminium or copper, and incorporate fins to increase surface area. Heat from a relatively small component is transferred to the larger heatsink; the equilibrium temperature of the component plus heatsink is much lower than the component's alone would be. Heat is carried away from the heatsink by convective or fan-forced airflow. Fan cooling is often used to cool processors and graphics cards that consume significant amounts of electrical energy. In a computer, a typical heat-generating component may be manufactured with a flat surface. A block of metal with a corresponding flat surface and finned construction, sometimes with an attached fan, is clamped to the component. To fill poorly conducting air gaps due to imperfectly flat and smooth surfaces, a thin layer of thermal grease, a thermal pad, or thermal adhesive may be placed between the component and heatsink.
Heat is removed from the heatsink by convection, to some extent by radiation, and possibly by conduction if the heatsink is in thermal contact with, say, the metal case. Inexpensive fan-cooled aluminium heatsinks are often used on standard desktop computers. Heatsinks with copper base-plates, or made of copper, have better thermal characteristics than those made of aluminium. A copper heatsink is more effective than an aluminium unit of the same size, which is relevant with regard to the high-power-consumption components used in high-performance computers.
Passive heatsinks are commonly found on older CPUs, parts that do not dissipate much power (such as the chipset), computers with low-power processors, and equipment where silent operation is critical and fan noise unacceptable.
Usually a heatsink is clamped to the integrated heat spreader (IHS), a flat metal plate the size of the CPU package which is part of the CPU assembly and spreads the heat locally. A thin layer of thermal compound is placed between them to compensate for surface imperfections. The spreader's primary purpose is to redistribute heat. The heatsink fins improve its efficiency.
Several brands of DDR2, DDR3, DDR4 and DDR5 DRAM memory modules are fitted with a finned heatsink clipped onto the top edge of the module. The same technique is used for video cards that use a finned passive heatsink on the GPU.
Higher-end M.2 SSDs can be prone to significant heat generation, and as a result these may be sold with a heatsink included, or alternatively a third-party heatsink may be attached by the user during installation.
Fan-cooled aluminium heatsinks were originally the norm for desktop computers, but nowadays many heatsinks feature copper base-plate, copper base-circle, or are entirely made of copper.
Dust tends to build up in the crevices of finned heatsinks, particularly with the high airflow produced by fans. This keeps the air away from the hot component, reducing cooling effectiveness; however, removing the dust restores effectiveness.
Peltier (thermoelectric) cooling
Peltier junctions are generally only around 10–15% as efficient as the ideal refrigerator (Carnot cycle), compared with 40–60% achieved by conventional compression cycle systems (reverse Rankine systems using compression/expansion). Due to this lower efficiency, thermoelectric cooling is generally only used in environments where the solid state nature (no moving parts, low maintenance, compact size, and orientation insensitivity) outweighs pure efficiency.
Modern TECs use several stacked units each composed of dozens or hundreds of thermocouples laid out next to each other, which allows for a substantial amount of heat transfer. A combination of bismuth and tellurium is most commonly used for the thermocouples.
As active heat pumps which consume power, TECs can produce temperatures below ambient, impossible with passive heatsinks, radiator-cooled liquid cooling, and heatpipe HSFs. However, while pumping heat, a Peltier module will typically consume more electric power than the heat amount being pumped.
It is also possible to use a Peltier element together with a high pressure refrigerant (two phase cooling) to cool the CPU.
Liquid cooling
Liquid cooling is a highly effective method of removing excess heat, with the most common heat transfer fluid in desktop PCs being (distilled) water. The advantages of water cooling over air cooling include water's higher specific heat capacity and thermal conductivity.
The principle used in a typical (active) liquid cooling system for computers is identical to that used in an automobile's internal combustion engine, with the water being circulated by a water pump through a water block mounted on the CPU (and sometimes additional components as GPU and northbridge) and out to a heat exchanger, typically a radiator. The radiator is itself usually cooled additionally by means of a fan. Besides a fan, it could possibly also be cooled by other means, such as a Peltier cooler (although Peltier elements are most commonly placed directly on top of the hardware to be cooled, and the coolant is used to conduct the heat away from the hot side of the Peltier element). A coolant reservoir is often also connected to the system.
Besides active liquid cooling systems, passive liquid cooling systems are also sometimes used. These systems often leave out a fan or a water pump, theoretically increasing their reliability and making them quieter than active systems. The downsides of these systems are that they are much less efficient in discarding the heat and thus also need to have much more coolantand thus a much bigger coolant reservoirgiving the coolant more time to cool down.
Liquids allow the transfer of more heat from the parts being cooled than air, making liquid cooling suitable for overclocking and high performance computer applications. Compared to air cooling, liquid cooling is also influenced less by the ambient temperature. Liquid cooling's comparatively low noise level compares favorably to that of air cooling, which can become quite noisy.
Disadvantages of liquid cooling include complexity and the potential for a coolant leak. Leaking water (and any additives in the water) can damage electronic components with which it comes into contact, and the need to test for and repair leaks makes for more complex and less reliable installations. (The first major foray into the field of liquid-cooled personal computers for general use, the high-end versions of Apple's Power Mac G5, was ultimately doomed by a propensity for coolant leaks.) An air-cooled heatsink is generally much simpler to build, install, and maintain than a water cooling solution, although CPU-specific water cooling kits can also be found, which may be just as easy to install as an air cooler. These are not limited to CPUs, and liquid cooling of GPU cards is also possible.
While originally limited to mainframe computers, liquid cooling has become a practice largely associated with overclocking in the form of either manufactured all-in-one (AIO) kits or do-it-yourself setups assembled from individually gathered parts. The past few years have seen an increase in the popularity of liquid cooling in pre-assembled, moderate to high performance, desktop computers. Sealed ("closed-loop") systems incorporating a small pre-filled radiator, fan, and waterblock simplify the installation and maintenance of water cooling at a slight cost in cooling effectiveness relative to larger and more complex setups. Liquid cooling is typically combined with air cooling, using liquid cooling for the hottest components, such as CPUs or GPUs, while retaining the simpler and cheaper air cooling for less demanding components.
The IBM Aquasar system uses hot water cooling to achieve energy efficiency, the water being used to heat buildings as well.
Since 2011, the effectiveness of water cooling has prompted a series of all-in-one (AIO) water cooling solutions. AIO solutions result in a much simpler to install the unit, and most units have been reviewed positively by review sites.
Heat pipes and vapor chambers
A heat pipe is a hollow tube containing a heat transfer liquid. The liquid absorbs heat and evaporates at one end of the pipe. The vapor travels to the other (cooler) end of the tube, where it condenses, giving up its latent heat. The liquid returns to the hot end of the tube by gravity or capillary action and repeats the cycle. Heat pipes have a much higher effective thermal conductivity than solid materials. For use in computers, the heatsink on the CPU is attached to a larger radiator heatsink. Both heatsinks are hollow, as is the attachment between them, creating one large heat pipe that transfers heat from the CPU to the radiator, which is then cooled using some conventional method. This method is usually used when space is tight, as in small form-factor PCs and laptops, or where no fan noise can be tolerated, as in audio production. Because of the efficiency of this method of cooling, many desktop CPUs and GPUs, as well as high end chipsets, use heat pipes or vapor chambers in addition to active fan-based cooling and passive heatsinks to remain within safe operating temperatures. A vapor chamber operates on the same principles as a heat pipe but takes on the form of a slab or sheet instead of a pipe. Heat pipes may be placed vertically on top and form part of vapor chambers. Vapor chambers may also be used on high-end smartphones.
Electrostatic air movement and corona discharge effect cooling
The cooling technology under development by Kronos and Thorn Micro Technologies employs a device called an ionic wind pump (also known as an electrostatic fluid accelerator). The basic operating principle of an ionic wind pump is corona discharge, an electrical discharge near a charged conductor caused by the ionization of the surrounding air.
The corona discharge cooler developed by Kronos works in the following manner: A high electric field is created at the tip of the cathode, which is placed on one side of the CPU. The high energy potential causes the oxygen and nitrogen molecules in the air to become ionized (positively charged) and create a corona (a halo of charged particles). Placing a grounded anode at the opposite end of the CPU causes the charged ions in the corona to accelerate towards the anode, colliding with neutral air molecules on the way. During these collisions, momentum is transferred from the ionized gas to the neutral air molecules, resulting in movement of gas towards the anode.
The advantages of the corona-based cooler are its lack of moving parts, thereby eliminating certain reliability issues and operating with a near-zero noise level and moderate energy consumption.
Soft cooling
Soft cooling is the practice of utilizing software to take advantage of CPU power saving technologies to minimize energy use. This is done using halt instructions to turn off or put in standby state CPU subparts that aren't being used or by underclocking the CPU. While resulting in lower total speeds, this can be very useful if overclocking a CPU to improve user experience rather than increase raw processing power, since it can prevent the need for noisier cooling. Contrary to what the term suggests, it is not a form of cooling but of reducing heat creation.
Undervolting
Undervolting is a practice of running the CPU or any other component with voltages below the device specifications. An undervolted component draws less power and thus produces less heat. The ability to do this varies by manufacturer, product line, and even different production runs of the same product (as well as that of other components in the system), but processors are often specified to use voltages higher than strictly necessary. This tolerance ensures that the processor will have a higher chance of performing correctly under sub-optimal conditions, such as a lower-quality motherboard or low power supply voltages. Below a certain limit, the processor will not function correctly, although undervolting too far does not typically lead to permanent hardware damage (unlike overvolting).
Undervolting is used for quiet systems, as less cooling is needed because of the reduction of heat production, allowing noisy fans to be omitted. It is also used when battery charge life must be maximized.
Chip-integrated
Conventional cooling techniques all attach their "cooling" component to the outside of the computer chip package. This "attaching" technique will always exhibit some thermal resistance, reducing its effectiveness. The heat can be more efficiently and quickly removed by directly cooling the local hot spots of the chip, within the package. At these locations, power dissipation of over 300 W/cm2 (typical CPU is less than 100 W/cm2) can occur, although future systems are expected to exceed 1000 W/cm2. This form of local cooling is essential to developing high power density chips. This ideology has led to the investigation of integrating cooling elements into the computer chip. Currently there are two techniques: micro-channel heatsinks, and jet impingement cooling.
In micro-channel heatsinks, channels are fabricated into the silicon chip (CPU), and coolant is pumped through them. The channels are designed with very large surface area which results in large heat transfers. Heat dissipation of 3000 W/cm2 has been reported with this technique. The heat dissipation can be further increased if two-phase flow cooling is applied. Unfortunately, the system requires large pressure drops, due to the small channels, and the heat flux is lower with dielectric coolants used in electronic cooling.
Another local chip cooling technique is jet impingement cooling. In this technique, a coolant is flowed through a small orifice to form a jet. The jet is directed toward the surface of the CPU chip, and can effectively remove large heat fluxes. Heat dissipation of over 1000 W/cm2 has been reported. The system can be operated at lower pressure in comparison to the micro-channel method. The heat transfer can be further increased using two-phase flow cooling and by integrating return flow channels (hybrid between micro-channel heatsinks and jet impingement cooling).
Phase-change cooling
Phase-change cooling is an extremely effective way to cool the processor. A vapor compression phase-change cooler is a unit that usually sits underneath the PC, with a tube leading to the processor. Inside the unit is a compressor of the same type as in an air conditioner. The compressor compresses a gas (or mixture of gases) which comes from the evaporator (CPU cooler discussed below). Then, the very hot high-pressure vapor is pushed into the condenser (heat dissipation device) where it condenses from a hot gas into a liquid, typically subcooled at the exit of the condenser then the liquid is fed to an expansion device (restriction in the system) to cause a drop in pressure a vaporize the fluid (cause it to reach a pressure where it can boil at the desired temperature); the expansion device used can be a simple capillary tube to a more elaborate thermal expansion valve. The liquid evaporates (changing phase), absorbing the heat from the processor as it draws extra energy from its environment to accommodate this change (see latent heat). The evaporation can produce temperatures reaching around . The liquid flows into the evaporator cooling the CPU, turning into a vapor at low pressure. At the end of the evaporator this gas flows down to the compressor and the cycle begins over again. This way, the processor can be cooled to temperatures ranging from , depending on the load, wattage of the processor, the refrigeration system (see refrigeration) and the gas mixture used. This type of system suffers from a number of issues (cost, weight, size, vibration, maintenance, cost of electricity, noise, need for a specialized computer tower) but, mainly, one must be concerned with dew point and the proper insulation of all sub-ambient surfaces that must be done (the pipes will sweat, dripping water on sensitive electronics).
Alternately, a new breed of the cooling system is being developed, inserting a pump into the thermosiphon loop. This adds another degree of flexibility for the design engineer, as the heat can now be effectively transported away from the heat source and either reclaimed or dissipated to ambient. Junction temperature can be tuned by adjusting the system pressure; higher pressure equals higher fluid saturation temperatures. This allows for smaller condensers, smaller fans, and/or the effective dissipation of heat in a high ambient temperature environment. These systems are, in essence, the next generation fluid cooling paradigm, as they are approximately 10 times more efficient than single-phase water. Since the system uses a dielectric as the heat transport medium, leaks do not cause a catastrophic failure of the electric system.
This type of cooling is seen as a more extreme way to cool components since the units are relatively expensive compared to the average desktop. They also generate a significant amount of noise, since they are essentially refrigerators; however, the compressor choice and air cooling system is the main determinant of this, allowing for flexibility for noise reduction based on the parts chosen.
A "thermosiphon" traditionally refers to a closed system consisting of several pipes and/or chambers, with a larger chamber containing a small reservoir of liquid (often having a boiling point just above ambient temperature, but not necessarily). The larger chamber is as close to the heat source and designed to conduct as much heat from it into the liquid as possible, for example, a CPU cold plate with the chamber inside it filled with the liquid. One or more pipes extend upward into some sort of radiator or similar heat dissipation area, and this is all set up such that the CPU heats the reservoir and liquid it contains, which begins boiling, and the vapor travels up the tube(s) into the radiator/heat dissipation area, and then after condensing, drips back down into the reservoir, or runs down the sides of the tube. This requires no moving parts, and is somewhat similar to a heat pump, except that capillary action is not used, making it potentially better in some sense (perhaps most importantly, better in that it is much easier to build, and much more customizable for specific use cases and the flow of coolant/vapor can be arranged in a much wider variety of positions and distances, and have far greater thermal mass and maximum capacity compared to heat pipes which are limited by the amount of coolant present and the speed and flow rate of coolant that capillary action can achieve with the wicking used, often sintered copper powder on the walls of the tube, which have a limited flow rate and capacity.)
Liquid nitrogen
As liquid nitrogen boils at , far below the freezing point of water, it is valuable as an extreme coolant for short overclocking sessions.
In a typical installation of liquid nitrogen cooling, a copper or aluminium pipe is mounted on top of the processor or graphics card. After the system has been heavily insulated against condensation, the liquid nitrogen is poured into the pipe, resulting in temperatures well below .
Evaporation devices ranging from cut out heatsinks with pipes attached to custom milled copper containers are used to hold the nitrogen as well as to prevent large temperature changes. However, after the nitrogen evaporates, it has to be refilled. In the realm of personal computers, this method of cooling is seldom used in contexts other than overclocking trial-runs and record-setting attempts, as the CPU will usually expire within a relatively short period of time due to temperature stress caused by changes in internal temperature.
Although liquid nitrogen is non-flammable, it can condense oxygen directly from air. Mixtures of liquid oxygen and flammable materials can be dangerously explosive.
Liquid nitrogen cooling is, generally, only used for processor benchmarking, due to the fact that continuous usage may cause permanent damage to one or more parts of the computer and, if handled in a careless way, can even harm the user, causing frostbite.
Liquid helium
Liquid helium, colder than liquid nitrogen, has also been used for cooling. Liquid helium boils at , and temperatures ranging from have been measured from the heatsink. However, liquid helium is more expensive and more difficult to store and use than liquid nitrogen. Also, extremely low temperatures can cause integrated circuits to stop functioning. Silicon-based semiconductors, for example, will freeze out at around .
Optimization
Cooling can be improved by several techniques which may involve additional expense or effort. These techniques are often used, in particular, by those who run parts of their computer (such as the CPU and GPU) at higher voltages and frequencies than specified by manufacturer (overclocking), which increases heat generation.
The installation of higher performance, non-stock cooling may also be considered modding. Many overclockers simply buy more efficient, and often, more expensive fan and heatsink combinations, while others resort to more exotic ways of computer cooling, such as liquid cooling, Peltier effect heatpumps, heat pipe or phase change cooling.
There are also some related practices that have a positive impact in reducing system temperatures:
Thermally conductive compounds
Often called Thermal Interface Material (TIM).
Perfectly flat surfaces in contact give optimal cooling, but perfect flatness and absence of microscopic air gaps is not practically possible, particularly in mass-produced equipment. A very thin skim of thermal compound, which is much more thermally conductive than air, though much less so than metal, can improve thermal contact and cooling by filling in the air gaps. If only a small amount of compound just sufficient to fill the gaps is used, the best temperature reduction will be obtained.
There is much debate about the merits of compounds, and overclockers often consider some compounds to be superior to others. The main consideration is to use the minimal amount of thermal compound required to even out surfaces, as the thermal conductivity of compound is typically 1/3 to 1/400 that of metal, though much better than air. The conductivity of the heatsink compound ranges from about 0.5 to 80W/mK (see articles); that of aluminium is about 200, that of air about 0.02. Heat-conductive pads are also used, often fitted by manufacturers to heatsinks. They are less effective than properly applied thermal compound, but simpler to apply and, if fixed to the heatsink, cannot be omitted by users unaware of the importance of good thermal contact, or replaced by a thick and ineffective layer of compound.
Unlike some techniques discussed here, the use of thermal compound or padding is almost universal when dissipating significant amounts of heat.
Heat sink lapping
Mass-produced CPU heat spreaders and heatsink bases are never perfectly flat or smooth; if these surfaces are placed in the best contact possible, there will be air gaps which reduce heat conduction. This can easily be mitigated by the use of thermal compound, but for the best possible results surfaces must be as flat as possible. This can be achieved by a laborious process known as lapping, which can reduce CPU temperature by typically .
Rounded cables
Most older PCs use flat ribbon cables to connect storage drives (IDE or SCSI). These large flat cables greatly impede airflow by causing drag and turbulence. Overclockers and modders often replace these with rounded cables, with the conductive wires bunched together tightly to reduce surface area. Theoretically, the parallel strands of conductors in a ribbon cable serve to reduce crosstalk (signal carrying conductors inducing signals in nearby conductors), but there is no empirical evidence of rounding cables reducing performance. This may be because the length of the cable is short enough so that the effect of crosstalk is negligible. Problems usually arise when the cable is not electromagnetically protected and the length is considerable, a more frequent occurrence with older network cables.
These computer cables can then be cable tied to the chassis or other cables to further increase airflow.
This is less of a problem with new computers that use serial ATA which has a much narrower cable.
Airflow
The colder the cooling medium (the air), the more effective the cooling. Cooling air temperature can be improved with these guidelines:
Supply cool air to the hot components as directly as possible. Examples are air snorkels and tunnels that feed outside air directly and exclusively to the CPU or GPU cooler. For example, the BTX case design prescribes a CPU air tunnel.
Expel warm air as directly as possible. Examples are: Conventional PC (ATX) power supplies blow the warm air out the back of the case. Many dual-slot graphics card designs blow the warm air through the cover of the adjacent slot. There are also some aftermarket coolers that do this. Some CPU cooling designs blow the warm air directly towards the back of the case, where it can be ejected by a case fan.
Air that has already been used to spot-cool a component should not be reused to spot-cool a different component (this follows from the previous items). The BTX case design violates this rule, since it uses the CPU cooler's exhaust to cool the chipset and often the graphics card. One may come across old or ultra-low-budget ATX cases which feature a PSU mount in the top. Most modern ATX cases do however have a PSU mount in the bottom of the case with a filtered air vent directly beneath the PSU.
Prefer cool intake air, avoid inhaling exhaust air (outside air above or near the exhausts). For example, a CPU cooling air duct at the back of a tower case would inhale warm air from a graphics card exhaust. Moving all exhausts to one side of the case, conventionally the back/top, helps to keep the intake air cool.
Hiding cables behind motherboard tray or simply apply ziptie and tucking cables away to provide unhindered airflow.
Fewer fans strategically placed will improve the airflow internally within the PC and thus lower the overall internal case temperature in relation to ambient conditions. The use of larger fans also improves efficiency and lowers the amount of waste heat along with the amount of noise generated by the fans while in operation.
There is little agreement on the effectiveness of different fan placement configurations, and little in the way of systematic testing has been done. For a rectangular PC (ATX) case, a fan in the front with a fan in the rear and one in the top has been found to be a suitable configuration. However, AMD's (somewhat outdated) system cooling guidelines notes that "A front cooling fan does not seem to be essential. In fact, in some extreme situations, testing showed these fans to be recirculating hot air rather than introducing cool air." It may be that fans in the side panels could have a similar detrimental effect – possibly through disrupting the normal air flow through the case. However, this is unconfirmed and probably varies with the configuration.
Air pressure
Loosely speaking, positive pressure means intake into the case is stronger than exhaust from the case. This configuration results in pressure inside of the case being higher than in its environment. Negative pressure means exhaust is stronger than intake. This results in internal air pressure being lower than in the environment. Both configurations have benefits and drawbacks, with positive pressure being the more popular of the two configurations. Negative pressure results in the case pulling air through holes and vents separate from the fans, as the internal gases will attempt to reach an equilibrium pressure with the environment. Consequently, this results in dust entering the computer in all locations. Positive pressure in combination with filtered intake solves this issue, as air will only incline to be exhausted through these holes and vents in order to reach an equilibrium with its environment. Dust is then unable to enter the case except through the intake fans, which need to possess dust filters.
Computer types
Desktops
Desktop computers typically use one or more fans for cooling. While almost all desktop power supplies have at least one built-in fan, power supplies should never draw heated air from within the case, as this results in higher PSU operating temperatures which decrease the PSU's energy efficiency, reliability and overall ability to provide a steady supply of power to the computer's internal components. For this reason, all modern ATX cases (with some exceptions found in ultra-low-budget cases) feature a power supply mount in the bottom, with a dedicated PSU air intake (often with its own filter) beneath the mounting location, allowing the PSU to draw cool air from beneath the case.
Most manufacturers recommend bringing cool, fresh air in at the bottom front of the case, and exhausting warm air from the top rear. If fans are fitted to force air into the case more effectively than it is removed, the pressure inside becomes higher than outside, referred to as a "positive" airflow (the opposite case is called "negative" airflow). Worth noting is that positive internal pressure only prevents dust accumulating in the case if the air intakes are equipped with dust filters. A case with negative internal pressure will suffer a higher rate of dust accumulation even if the intakes are filtered, as the negative pressure will draw dust in through any available opening in the case
The air flow inside the typical desktop case is usually not strong enough for a passive CPU heatsink. Most desktop heatsinks are active including one or even multiple directly attached fans or blowers.
Servers
Server coolers
Each server can have an independent internal cooler system; Server cooling fans in (1 U) enclosures are usually located in the middle of the enclosure, between the hard drives at the front and passive CPU heatsinks at the rear. Larger (higher) enclosures also have exhaust fans, and from approximately 4U they may have active heatsinks. Power supplies generally have their own rear-facing exhaust fans.
Rack-mounted coolers
Rack cabinet is a typical enclosure for horizontally mounted servers. Air typically drawn in at the front of the rack and exhausted at the rear. Each cabinet can have additional cooling options; for example, they can have a Close Coupled Cooling attachable module or integrated with cabinet elements (like cooling doors in iDataPlex server rack).
Another way of accommodating large numbers of systems in a small space is to use blade chassis, oriented vertically rather than horizontally, to facilitate convection. Air heated by the hot components tends to rise, creating a natural air flow along the boards (stack effect), cooling them. Some manufacturers take advantage of this effect.
Data center cooling
Because data centers typically contain large numbers of computers and other power-dissipating devices, they risk equipment overheating; extensive HVAC systems are used to prevent this. Often a raised floor is used so the area under the floor may be used as a large plenum for cooled air from a CRAC (Computer Room Air Conditioner) or a CRAH (Computer Room Air Handler) and power cabling. A plenum made with a false ceiling can also be present. Hot Aisle containment or cold aisle containment are also used in datacenters to improve cooling efficiency. Alternatively slab floors can be used which are similar to conventional floors, and overhead ducts can be used for cooling.
Direct Contact Liquid Cooling has emerged more efficient than air cooling options, resulting in smaller footprint, lower capital requirements and lower operational costs than air cooling. It uses warm liquid instead of air to move heat away from the hottest components. Energy efficiency gains from liquid cooling is also driving its adoption. Single and dual/two phase immersion/open tub cooling and single and dual phase direct-to-chip cooling as well as immersion cooling confined to individual server blades have also been proposed for use in data centers. In-row cooling, rack cooling, rear door heat exchangers, racktop cooling which places heat exchangers above the rack, overhead cooling above aisles or fan walls/thermal walls in a data center can also be used. Direct Liquid Cooling (DLC) with cold plates for cooling chips in servers can be used due to the higher heat removal capacities of these systems. These systems can either cool some or all components on a server, using rubber or copper tubing respectively. Rear door heat exchangers were traditionally used for cooling high heat densities in data centers, but these did not see widespread adoption. They can be cooled with refrigerant or chilled water. Those cooled with chilled water can either be active, which have fans, or passive, which have no fans. Liquid to air heat exchangers (radiators) can be used to cool servers cooled with direct to chip liquid cooling, in order to avoid installation of facility water piping. These heat exchangers can be installed separately from racks, or as a rear door on a rack.
Laptops
Laptops present a difficult mechanical airflow design, power dissipation, and cooling challenge. Constraints specific to laptops include: the device as a whole has to be as light as possible; the form factor has to be built around the standard keyboard layout; users are very close, so noise must be kept to a minimum, and the case exterior temperature must be kept low enough to be used on a lap. Cooling generally uses forced air cooling but heat pipes and the use of the metal chassis or case as a passive heatsink are also common. Solutions to reduce heat include using lower power-consumption ARM or Intel Atom processors.
Mobile devices
Mobile devices such as phones usually have no discrete cooling systems, as mobile CPU and GPU chips are designed for maximum power efficiency due to the constraints of the device's battery. Some higher performance devices may include a heat spreader that aids in transferring heat to the external case of a phone or tablet.
| Technology | Computer hardware | null |
799405 | https://en.wikipedia.org/wiki/Shear%20mapping | Shear mapping | In plane geometry, a shear mapping is an affine transformation that displaces each point in a fixed direction by an amount proportional to its signed distance from a given line parallel to that direction. This type of mapping is also called shear transformation, transvection, or just shearing. The transformations can be applied with a shear matrix or transvection, an elementary matrix that represents the addition of a multiple of one row or column to another. Such a matrix may be derived by taking the identity matrix and replacing one of the zero elements with a non-zero value.
An example is the linear map that takes any point with coordinates to the point . In this case, the displacement is horizontal by a factor of 2 where the fixed line is the -axis, and the signed distance is the -coordinate. Note that points on opposite sides of the reference line are displaced in opposite directions.
Shear mappings must not be confused with rotations. Applying a shear map to a set of points of the plane will change all angles between them (except straight angles), and the length of any line segment that is not parallel to the direction of displacement. Therefore, it will usually distort the shape of a geometric figure, for example turning squares into parallelograms, and circles into ellipses. However a shearing does preserve the area of geometric figures and the alignment and relative distances of collinear points. A shear mapping is the main difference between the upright and slanted (or italic) styles of letters.
The same definition is used in three-dimensional geometry, except that the distance is measured from a fixed plane. A three-dimensional shearing transformation preserves the volume of solid figures, but changes areas of plane figures (except those that are parallel to the displacement).
This transformation is used to describe laminar flow of a fluid between plates, one moving in a plane above and parallel to the first.
In the general -dimensional Cartesian space the distance is measured from a fixed hyperplane parallel to the direction of displacement. This geometric transformation is a linear transformation of that preserves the -dimensional measure (hypervolume) of any set.
Definition
Horizontal and vertical shear of the plane
In the plane , a horizontal shear (or shear parallel to the -axis) is a function that takes a generic point with coordinates to the point ; where is a fixed parameter, called the shear factor.
The effect of this mapping is to displace every point horizontally by an amount proportionally to its -coordinate. Any point above the -axis is displaced to the right (increasing ) if , and to the left if . Points below the -axis move in the opposite direction, while points on the axis stay fixed.
Straight lines parallel to the -axis remain where they are, while all other lines are turned (by various angles) about the point where they cross the -axis. Vertical lines, in particular, become oblique lines with slope Therefore, the shear factor is the cotangent of the shear angle between the former verticals and the -axis. (In the example on the right the square is tilted by 30°, so the shear angle is 60°.)
If the coordinates of a point are written as a column vector (a 2×1 matrix), the shear mapping can be written as multiplication by a 2×2 matrix:
A vertical shear (or shear parallel to the -axis) of lines is similar, except that the roles of and are swapped. It corresponds to multiplying the coordinate vector by the transposed matrix:
The vertical shear displaces points to the right of the -axis up or down, depending on the sign of . It leaves vertical lines invariant, but tilts all other lines about the point where they meet the -axis. Horizontal lines, in particular, get tilted by the shear angle to become lines with slope .
Composition
Two or more shear transformations can be combined.
If two shear matrices are and
then their composition matrix is
which also has determinant 1, so that area is preserved.
In particular, if , we have
which is a positive definite matrix.
Higher dimensions
A typical shear matrix is of the form
This matrix shears parallel to the axis in the direction of the fourth dimension of the underlying vector space.
A shear parallel to the axis results in and . In matrix form:
Similarly, a shear parallel to the axis has and . In matrix form:
In 3D space this matrix shear the YZ plane into the diagonal plane passing through these 3 points:
The determinant will always be 1, as no matter where the shear element is placed, it will be a member of a skew-diagonal that also contains zero elements (as all skew-diagonals have length at least two) hence its product will remain zero and will not contribute to the determinant. Thus every shear matrix has an inverse, and the inverse is simply a shear matrix with the shear element negated, representing a shear transformation in the opposite direction. In fact, this is part of an easily derived more general result: if is a shear matrix with shear element , then is a shear matrix whose shear element is simply . Hence, raising a shear matrix to a power multiplies its shear factor by .
Properties
If is an shear matrix, then:
has rank and therefore is invertible
1 is the only eigenvalue of , so and
the eigenspace of (associated with the eigenvalue 1) has dimensions.
is defective
is asymmetric
may be made into a block matrix by at most 1 column interchange and 1 row interchange operation
the area, volume, or any higher order interior capacity of a polytope is invariant under the shear transformation of the polytope's vertices.
General shear mappings
For a vector space and subspace , a shear fixing translates all vectors in a direction parallel to .
To be more precise, if is the direct sum of and , and we write vectors as
correspondingly, the typical shear fixing is
where is a linear mapping from into . Therefore in block matrix terms can be represented as
Applications
The following applications of shear mapping were noted by William Kingdon Clifford:
"A succession of shears will enable us to reduce any figure bounded by straight lines to a triangle of equal area."
"... we may shear any triangle into a right-angled triangle, and this will not alter its area. Thus the area of any triangle is half the area of the rectangle on the same base and with height equal to the perpendicular on the base from the opposite angle."
The area-preserving property of a shear mapping can be used for results involving area. For instance, the Pythagorean theorem has been illustrated with shear mapping as well as the related geometric mean theorem.
Shear matrices are often used in computer graphics.
An algorithm due to Alan W. Paeth uses a sequence of three shear mappings (horizontal, vertical, then horizontal again) to rotate a digital image by an arbitrary angle. The algorithm is very simple to implement, and very efficient, since each step processes only one column or one row of pixels at a time.
In typography, normal text transformed by a shear mapping results in oblique type.
In pre-Einsteinian Galilean relativity, transformations between frames of reference are shear mappings called Galilean transformations. These are also sometimes seen when describing moving reference frames relative to a "preferred" frame, sometimes referred to as absolute time and space.
| Mathematics | Geometry: General | null |
799436 | https://en.wikipedia.org/wiki/Pitchfork | Pitchfork | A pitchfork or hay fork is an agricultural tool used to pitch loose material, such as hay, straw, manure, or leaves. It has a long handle and usually two to five thin tines designed to efficiently move such materials.
The term is also applied colloquially, but inaccurately, to the garden fork. While similar in appearance, the garden fork is shorter and stockier than the pitchfork, with three or four thicker tines intended for turning or loosening the soil of gardens.
Alternative terms
In some parts of England, a pitchfork is known as a prong. In parts of Ireland, the term sprong is used to refer specifically to a four-pronged pitchfork.
Description
The typical pitchfork consists of a wooden shaft bearing two to five slightly curved metal tines fixed to one end of a handle. These are typically made of steel, wrought iron, or some other alloy, though historically wood or bamboo were used. Unlike a garden fork, a pitchfork lacks a grab at the end of its handle.
Pitchforks with few tines set far apart are typically used for bulky material such as hay or straw; those with more and more closely spaced are used for looser materials such as silage, manure, leaves, or compost.
History
In Europe, the pitchfork was first used in the Early Middle Ages, at about the same time as the harrow. These were made entirely of wood.
The pitchfork is occasionally employed as an improvised weapon, as in a mob or riot.
In popular culture
Artwork
Paintings by various artists depict a wide variety of pitchforks in use and at rest. A notable American work is American Gothic (1930) by Grant Wood, which features a three-pronged tool.
Politics
Because of its association with peasantry and farming, the pitchfork has been used as a populist symbol and appended as a nickname for certain leading populist figures, such as "Pitchfork" Ben Tillman and "Pitchfork" Pat Buchanan.
The Gangster Disciples, a street gang in the Midwestern United States, use a three-pointed pitchfork as one of their symbols.
Venezuelan far-right political party, New Order use three-pointed pitchforks as their symbol.
Religious symbolism
The pitchfork is often used in lieu of the visually similar weapon, the trident, in popular portrayals and satire of Christian demonology. Many humorous cartoons, both animated and otherwise, feature a caricature of a demon ostensibly wielding a "pitchfork" (often actually a trident) sitting on one shoulder of a protagonist, opposite an angel on the other.
The Hellenistic deity Hades wields a bident, a two-pronged weapon similar in form to a pitchfork but actually related to the trident in design and purpose.
| Technology | Agricultural tools | null |
799876 | https://en.wikipedia.org/wiki/Electric%20susceptibility | Electric susceptibility | In electricity (electromagnetism), the electric susceptibility (; Latin: susceptibilis "receptive") is a dimensionless proportionality constant that indicates the degree of polarization of a dielectric material in response to an applied electric field. The greater the electric susceptibility, the greater the ability of a material to polarize in response to the field, and thereby reduce the total electric field inside the material (and store energy). It is in this way that the electric susceptibility influences the electric permittivity of the material and thus influences many other phenomena in that medium, from the capacitance of capacitors to the speed of light.
Definition for linear dielectrics
If a dielectric material is a linear dielectric, then electric susceptibility is defined as the constant of proportionality (which may be a tensor) relating an electric field E to the induced dielectric polarization density P such that
where
is the polarization density;
is the electric permittivity of free space (electric constant);
is the electric susceptibility;
is the electric field.
In materials where susceptibility is anisotropic (different depending on direction), susceptibility is represented as a tensor known as the susceptibility tensor. Many linear dielectrics are isotropic, but it is possible nevertheless for a material to display behavior that is both linear and anisotropic, or for a material to be non-linear but isotropic. Anisotropic but linear susceptibility is common in many crystals.
The susceptibility is related to its relative permittivity (dielectric constant) by
so in the case of a vacuum,
At the same time, the electric displacement D is related to the polarization density P by the following relation:
where
Molecular polarizability
A similar parameter exists to relate the magnitude of the induced dipole moment p of an individual molecule to the local electric field E that induced the dipole. This parameter is the molecular polarizability (α), and the dipole moment resulting from the local electric field Elocal is given by:
This introduces a complication however, as locally the field can differ significantly from the overall applied field. We have:
where P is the polarization per unit volume, and N is the number of molecules per unit volume contributing to the polarization. Thus, if the local electric field is parallel to the ambient electric field, we have:
Thus only if the local field equals the ambient field can we write:
Otherwise, one should find a relation between the local and the macroscopic field. In some materials, the Clausius–Mossotti relation holds and reads
Ambiguity in the definition
The definition of the molecular polarizability depends on the author. In the above definition,
and are in SI units and the molecular polarizability has the dimension of a volume (m3). Another definition would be to keep SI units and to integrate into :
In this second definition, the polarizability would have the SI unit of C.m2/V. Yet another definition exists where and are expressed in the cgs system and is still defined as
Using the cgs units gives the dimension of a volume, as in the first definition, but with a value that is lower.
Nonlinear susceptibility
In many materials the polarizability starts to saturate at high values of electric field. This saturation can be modelled by a nonlinear susceptibility. These susceptibilities are important in nonlinear optics and lead to effects such as second-harmonic generation (such as used to convert infrared light into visible light, in green laser pointers).
The standard definition of nonlinear susceptibilities in SI units is via a Taylor expansion of the polarization's reaction to electric field:
(Except in ferroelectric materials, the built-in polarization is zero, .)
The first susceptibility term, , corresponds to the linear susceptibility described above. While this first term is dimensionless, the subsequent nonlinear susceptibilities have units of .
The nonlinear susceptibilities can be generalized to anisotropic materials in which the susceptibility is not uniform in every direction. In these materials, each susceptibility becomes an ()-degree tensor.
Dispersion and causality
In general, a material cannot polarize instantaneously in response to an applied field, and so the more general formulation as a function of time is
That is, the polarization is a convolution of the electric field at previous times with time-dependent susceptibility given by . The upper limit of this integral can be extended to infinity as well if one defines for . An instantaneous response corresponds to Dirac delta function susceptibility .
It is more convenient in a linear system to take the Fourier transform and write this relationship as a function of frequency. Due to the convolution theorem, the integral becomes a product,
This has a similar form to the Clausius–Mossotti relation:
This frequency dependence of the susceptibility leads to frequency dependence of the permittivity. The shape of the susceptibility with respect to frequency characterizes the dispersion properties of the material.
Moreover, the fact that the polarization can only depend on the electric field at previous times (i.e. for ), a consequence of causality, imposes Kramers–Kronig constraints on the susceptibility .
| Physical sciences | Electrostatics | Physics |
799986 | https://en.wikipedia.org/wiki/Fraunhofer%20diffraction | Fraunhofer diffraction | In optics, the Fraunhofer diffraction equation is used to model the diffraction of waves when plane waves are incident on a diffracting object, and the diffraction pattern is viewed at a sufficiently long distance (a distance satisfying Fraunhofer condition) from the object (in the far-field region), and also when it is viewed at the focal plane of an imaging lens. In contrast, the diffraction pattern created near the diffracting object and (in the near field region) is given by the Fresnel diffraction equation.
The equation was named in honor of Joseph von Fraunhofer although he was not actually involved in the development of the theory.
This article explains where the Fraunhofer equation can be applied, and shows Fraunhofer diffraction patterns for various apertures. A detailed mathematical treatment of Fraunhofer diffraction is given in Fraunhofer diffraction equation.
Equation
When a beam of light is partly blocked by an obstacle, some of the light is scattered around the object, light and dark bands are often seen at the edge of the shadow – this effect is known as diffraction. These effects can be modelled using the Huygens–Fresnel principle; Huygens postulated that every point on a wavefront acts as a source of spherical secondary wavelets and the sum of these secondary wavelets determines the form of the proceeding wave at any subsequent time, while Fresnel developed an equation using the Huygens wavelets together with the principle of superposition of waves, which models these diffraction effects quite well.
It is generally not straightforward to calculate the wave amplitude given by the sum of the secondary wavelets (The wave sum is also a wave.), each of which has its own amplitude, phase, and oscillation direction (polarization), since this involves addition of many waves of varying amplitude, phase, and polarization. When two light waves as electromagnetic fields are added together (vector sum), the amplitude of the wave sum depends on the amplitudes, the phases, and even the polarizations of individual waves. On a certain direction where electromagnetic wave fields are projected (or considering a situation where two waves have the same polarization), two waves of equal (projected) amplitude which are in phase (same phase) give the amplitude of the resultant wave sum as double the individual wave amplitudes, while two waves of equal amplitude which are in opposite phases give the zero amplitude of the resultant wave as they cancel out each other. Generally, a two-dimensional integral over complex variables has to be solved and in many cases, an analytic solution is not available.
The Fraunhofer diffraction equation is a simplified version of Kirchhoff's diffraction formula and it can be used to model light diffraction when both a light source and a viewing plane (a plane of observation where the diffracted wave is observed) are effectively infinitely distant from a diffracting aperture. With a sufficiently distant light source from a diffracting aperture, the incident light to the aperture is effectively a plane wave so that the phase of the light at each point on the aperture is the same. At a sufficiently distant plane of observation from the aperture, the phase of the wave coming from each point on the aperture varies linearly with the point position on the aperture, making the calculation of the sum of the waves at an observation point on the plane of observation relatively straightforward in many cases. Even the amplitudes of the secondary waves coming from the aperture at the observation point can be treated as same or constant for a simple diffraction wave calculation in this case. Diffraction in such a geometrical requirement is called Fraunhofer diffraction, and the condition where Fraunhofer diffraction is valid is called Fraunhofer condition, as shown in the right box. A diffracted wave is often called Far field if it at least partially satisfies Fraunhofer condition such that the distance between the aperture and the observation plane is .
For example, if a 0.5 mm diameter circular hole is illuminated by a laser light with 0.6 μm wavelength, then Fraunhofer diffraction occurs if the viewing distance is greater than 1000 mm.
Derivation of Fraunhofer condition
The derivation of Fraunhofer condition here is based on the geometry described in the right box. The diffracted wave path r2 can be expressed in terms of another diffracted wave path r1 and the distance b between two diffracting points by using the law of cosines;
This can be expanded by calculating the expression's Taylor series to second order with respect to ,
The phase difference between waves propagating along the paths r2 and r1 are, with the wavenumber where λ is the light wavelength,
If so , then the phase difference is . The geometrical implication from this expression is that the paths r2 and r1 are approximately parallel with each other. Since there can be a diffraction - observation plane, the diffracted wave path whose angle with respect to a straight line parallel to the optical axis is close to 0, this approximation condition can be further simplified as where L is the distance between two planes along the optical axis. Due to the fact that an incident wave on a diffracting plane is effectively a plane wave if where L is the distance between the diffracting plane and the point wave source is satisfied, Fraunhofer condition is where L is the smaller of the two distances, one is between the diffracting plane and the plane of observation and the other is between the diffracting plane and the point wave source.
Focal plane of a positive lens as the far field plane
In the far field, propagation paths for wavelets from every point on an aperture to a point of observation are approximately parallel, and a positive lens (focusing lens) focuses parallel rays toward the lens to a point on the focal plane (the focus point position on the focal plane depends on the angle of the parallel rays with respect to the optical axis). So, if a positive lens with a sufficiently long focal length (so that differences between electric field orientations for wavelets can be ignored at the focus) is placed after an aperture, then the lens practically makes the Fraunhofer diffraction pattern of the aperture on its focal plane as the parallel rays meet each other at the focus.
Examples
In each of these examples, the aperture is illuminated by a monochromatic plane wave at normal incidence.
Diffraction by a narrow rectangular slit
The width of the slit is . The Fraunhofer diffraction pattern is shown in the image together with a plot of the intensity vs. angle . The pattern has maximum intensity at , and a series of peaks of decreasing intensity. Most of the diffracted light falls between the first minima. The angle, , subtended by these two minima is given by:
Thus, the smaller the aperture, the larger the angle subtended by the diffraction bands. The size of the central band at a distance is given by
For example, when a slit of width 0.5 mm is illuminated by light of wavelength 0.6 μm, and viewed at a distance of 1000 mm, the width of the central band in the diffraction pattern is 2.4 mm.
The fringes extend to infinity in the direction since the slit and illumination also extend to infinity.
If , the intensity of the diffracted light does not fall to zero, and if , the diffracted wave is cylindrical.
Semi-quantitative analysis of single-slit diffraction
We can find the angle at which a first minimum is obtained in the diffracted light by the following reasoning. Consider the light diffracted at an angle where the distance is equal to the wavelength of the illuminating light. The width of the slit is the distance . The component of the wavelet emitted from the point A which is travelling in the direction is in anti-phase with the wave from the point at middle of the slit, so that the net contribution at the angle from these two waves is zero. The same applies to the points just below and , and so on. Therefore, the amplitude of the total wave travelling in the direction is zero. We have:
The angle subtended by the first minima on either side of the centre is then, as above:
There is no such simple argument to enable us to find the maxima of the diffraction pattern.
Single-slit diffraction using Huygens' principle
We can develop an expression for the far field of a continuous array of point sources of uniform amplitude and of the same phase. Let the array of length a be parallel to the y axis with its center at the origin as indicated in the figure to the right. Then the differential field is:
where . However and integrating from to ,
where .
Integrating we then get
Letting where the array length in radians is , then,
Diffraction by a rectangular aperture
The form of the diffraction pattern given by a rectangular aperture is shown in the figure on the right (or above, in tablet format). There is a central semi-rectangular peak, with a series of horizontal and vertical fringes. The dimensions of the central band are related to the dimensions of the slit by the same relationship as for a single slit so that the larger dimension in the diffracted image corresponds to the smaller dimension in the slit. The spacing of the fringes is also inversely proportional to the slit dimension.
If the illuminating beam does not illuminate the whole vertical length of the slit, the spacing of the vertical fringes is determined by the dimensions of the illuminating beam. Close examination of the double-slit diffraction pattern below shows that there are very fine horizontal diffraction fringes above and below the main spot, as well as the more obvious horizontal fringes.
Diffraction by a circular aperture
The diffraction pattern given by a circular aperture is shown in the figure on the right. This is known as the Airy diffraction pattern. It can be seen that most of the light is in the central disk. The angle subtended by this disk, known as the Airy disk, is
where is the diameter of the aperture.
The Airy disk can be an important parameter in limiting the ability of an imaging system to resolve closely located objects.
Diffraction by an aperture with a Gaussian profile
The diffraction pattern obtained given by an aperture with a Gaussian profile, for example, a photographic slide whose transmissivity has a Gaussian variation is also a Gaussian function. The form of the function is plotted on the right (above, for a tablet), and it can be seen that, unlike the diffraction patterns produced by rectangular or circular apertures, it has no secondary rings. This technique can be used in a process called apodization—the aperture is covered by a Gaussian filter, giving a diffraction pattern with no secondary rings.
The output profile of a single mode laser beam may have a Gaussian intensity profile and the diffraction equation can be used to show that it maintains that profile however far away it propagates from the source.
Diffraction by a double slit
In the double-slit experiment, the two slits are illuminated by a single light beam. If the width of the slits is small enough (less than the wavelength of the light), the slits diffract the light into cylindrical waves. These two cylindrical wavefronts are superimposed, and the amplitude, and therefore the intensity, at any point in the combined wavefronts depends on both the magnitude and the phase of the two wavefronts. These fringes are often known as Young's fringes.
The angular spacing of the fringes is given by
The spacing of the fringes at a distance from the slits is given by
where is the separation of the slits.
The fringes in the picture were obtained using the yellow light from a sodium light (wavelength = 589 nm), with slits separated by 0.25 mm, and projected directly onto the image plane of a digital camera.
Double-slit interference fringes can be observed by cutting two slits in a piece of card, illuminating with a laser pointer, and observing the diffracted light at a distance of 1 m. If the slit separation is 0.5 mm, and the wavelength of the laser is 600 nm, then the spacing of the fringes viewed at a distance of 1 m would be 1.2 mm.
Semi-quantitative explanation of double-slit fringes
The difference in phase between the two waves is determined by the difference in the distance travelled by the two waves.
If the viewing distance is large compared with the separation of the slits (the far field), the phase difference can be found using the geometry shown in the figure. The path difference between two waves travelling at an angle is given by
When the two waves are in phase, i.e. the path difference is equal to an integral number of wavelengths, the summed amplitude, and therefore the summed intensity is maximal, and when they are in anti-phase, i.e. the path difference is equal to half a wavelength, one and a half wavelengths, etc., then the two waves cancel, and the summed intensity is zero. This effect is known as interference.
The interference fringe maxima occur at angles
where is the wavelength of the light. The angular spacing of the fringes is given by
When the distance between the slits and the viewing plane is , the spacing of the fringes is equal to and is the same as above:
Diffraction by a grating
A grating is defined in Born and Wolf as "any arrangement which imposes on an incident wave a periodic variation of amplitude or phase, or both".
A grating whose elements are separated by diffracts a normally incident beam of light into a set of beams, at angles given by:
This is known as the grating equation. The finer the grating spacing, the greater the angular separation of the diffracted beams.
If the light is incident at an angle , the grating equation is:
The detailed structure of the repeating pattern determines the form of the individual diffracted beams, as well as their relative intensity while the grating spacing always determines the angles of the diffracted beams.
The image on the right shows a laser beam diffracted by a grating into = 0, and ±1 beams. The angles of the first order beams are about 20°; if we assume the wavelength of the laser beam is 600 nm, we can infer that the grating spacing is about 1.8 μm.
Semi-quantitative explanation
A simple grating consists of a series of slits in a screen. If the light travelling at an angle from each slit has a path difference of one wavelength with respect to the adjacent slit, all these waves will add together, so that the maximum intensity of the diffracted light is obtained when:
This is the same relationship that is given above.
| Physical sciences | Waves | Physics |
800373 | https://en.wikipedia.org/wiki/Mosasaur | Mosasaur | Mosasaurs (from Latin Mosa meaning the 'Meuse', and Greek meaning 'lizard') are an extinct group of large aquatic reptiles within the family Mosasauridae that lived during the Late Cretaceous. Their first fossil remains were discovered in a limestone quarry at Maastricht on the Meuse in 1764. They belong to the order Squamata, which includes lizards and snakes.
During the last 20 million years of the Cretaceous period (Turonian–Maastrichtian ages), with the extinction of the ichthyosaurs and pliosaurs, mosasaurids became the dominant marine predators. They themselves became extinct as a result of the K-Pg event at the end of the Cretaceous period, about 66 million years ago.
Description
Mosasaurs breathed air, were powerful swimmers, and were well-adapted to living in the warm, shallow inland seas prevalent during the Late Cretaceous period. Mosasaurs were so well adapted to this environment that they most likely gave birth to live young, rather than returning to the shore to lay eggs as sea turtles do.
The smallest-known mosasaur was Dallasaurus turneri, which was less than long. Larger mosasaurs were more typical, with many species growing longer than . Mosasaurus hoffmannii, the largest known species reached up to , but it has been considered to be probably overestimated by Cleary et al. (2018). Currently, the largest publicly exhibited mosasaur skeleton in the world is on display at the Canadian Fossil Discovery Centre in Morden, Manitoba. The specimen, nicknamed "Bruce", is just over long, but this might be an overestimate as the skeleton was assembled for display prior to a 2010 reassessment of the species that found its original number of vertebrae to be exaggerated, implying that the actual size of the animal was likely smaller.
Mosasaurs had a body shape similar to that of modern-day monitor lizards (varanids), but were more elongated and streamlined for swimming. Their limb bones were reduced in length and their paddles were formed by webbing between their long finger and toe bones. Their tails were broad, and supplied their locomotive power.
Until recently, mosasaurs were assumed to have swum in a method similar to the one used today by conger eels and sea snakes, undulating their entire bodies from side to side. However, new evidence suggests that many advanced mosasaurs had large, crescent-shaped flukes on the ends of their tails, similar to those of sharks and some ichthyosaurs. Rather than use snake-like undulations, their bodies probably remained stiff to reduce drag through the water, while their tails provided strong propulsion. These animals may have lurked and pounced rapidly and powerfully on passing prey, rather than chasing after it. At least some species were also capable of aquaflight, flapping their flippers like sea lions.
Early reconstructions showed mosasaurs with dorsal crests running the length of their bodies, which were based on misidentified remains of tracheal cartilage. By the time this error was discovered, depicting mosasaurs with such crests in artwork had already become a trend.
Paleobiology
Mosasaurs had double-hinged jaws and flexible skulls (much like those of snakes), which enabled them to gulp down their prey almost whole. A skeleton of Tylosaurus proriger from South Dakota included remains of the diving seabird Hesperornis, a marine bony fish, a possible shark, and another, smaller mosasaur (Clidastes). Mosasaur bones have also been found with shark teeth embedded in them.
One of the food items of mosasaurs were ammonites, molluscs with shells similar to those of Nautilus, which were abundant in the Cretaceous seas. Holes have been found in fossil shells of some ammonites, mainly Pachydiscus and Placenticeras. These were once interpreted as a result of limpets attaching themselves to the ammonites, but the triangular shape of the holes, their size, and their presence on both sides of the shells, corresponding to upper and lower jaws, is evidence of the bite of medium-sized mosasaurs. Whether this behaviour was common across all size classes of mosasaurs is not clear.
Virtually all forms were active predators of fish and ammonites; a few, such as Globidens, had blunt, spherical teeth, specialized for crushing mollusk shells. The smaller genera, such as Platecarpus and Dallasaurus, which were about long, probably fed on fish and other small prey. The smaller mosasaurs may have spent some time in fresh water, hunting for food. The largest mosasaur Mosasaurus hoffmannii was the apex predator of the Late Cretaceous oceans, reaching more than in length and weighing up to in body mass.
Soft tissue
Despite the many mosasaur remains collected worldwide, knowledge of the nature of their skin coverings remains in its early stages. Few mosasaurid specimens collected from around the world retain fossilized scale imprints. This lack may be due to the delicate nature of the scales, which nearly eliminates the possibility of preservation, in addition to the preservation sediment types and the marine conditions under which the preservation occurred. Until the discovery of several mosasaur specimens with remarkably well-preserved scale imprints from late Maastrichtian deposits of the Muwaqqar Chalk Marl Formation of Harrana in Jordan, knowledge of the nature of mosasaur integument was mainly based on very few accounts describing early mosasaur fossils dating back to the upper Santonian–lower Campanian, such as the famous Tylosaurus specimen (KUVP-1075) from Gove County, Kansas.
Material from Jordan has shown that the bodies of mosasaurs, as well as the membranes between their fingers and toes, were covered with small, overlapping, diamond-shaped scales resembling those of snakes. Much like those of modern reptiles, mosasaur scales varied across the body in type and size. In Harrana specimens, two types of scales were observed on a single specimen: keeled scales covering the upper regions of the body and smooth scales covering the lower. As ambush predators, lurking and quickly capturing prey using stealth tactics, they may have benefited from the nonreflective, keeled scales. Additionally, mosasaurs had large pectoral girdles, and such genera as Plotosaurus may have used their front flippers in a breaststroke motion to gain added bursts of speed during an attack on prey.
More recently, a fossil of Platecarpus tympaniticus has been found that preserved not only skin impressions, but also internal organs. Several reddish areas in the fossil may represent the heart, lungs, and kidneys. The trachea is also preserved, along with part of what may be the retina in the eye. The placement of the kidneys is farther forward in the abdomen than it is in monitor lizards, and is more similar to those of cetaceans. As in cetaceans, the bronchi leading to the lungs run parallel to each other instead of splitting apart from one another as in monitors and other terrestrial reptiles. In mosasaurs, these features may be internal adaptations to fully marine lifestyles.
In 2011, collagen protein was recovered from a Prognathodon humerus dated to the Cretaceous.
In 2005, a case study by A.S. Schulp, E.W.A Mulder, and K. Schwenk outlined the fact that mosasaurs had paired fenestrae in their palates. In monitor lizards and snakes, paired fenestrae are associated with a forked tongue, which is flicked in and out to detect chemical traces and provide a directional sense of smell. They therefore proposed that mosasaurs probably also had a sensitive forked tongue.
Metabolism
A study published in 2016 by T. Lyn Harrell, Alberto Pérez-Huerta and Celina Suarez showed that mosasaurs were endothermic. The study contradicted findings published in 2010 indicating mosasaurs were ectothermic. The 2010 study did not use warm-blooded animals for comparison but analogous groups of common marine animals. Based on comparisons with modern warm-blooded animals and fossils of known cold-blooded animals from the same time period, the 2016 study found mosasaurs likely had body temperatures similar to those of contemporary seabirds and were able to internally regulate their temperatures to remain warmer than the surrounding water.
Coloration
The coloration of mosasaurs was unknown until 2014, when the findings of Johan Lindgren of Lund University and colleagues revealed the pigment melanin in the fossilized scales of a mosasaur. Mosasaurs were likely countershaded, with dark backs and light underbellies, much like a great white shark or leatherback sea turtle, the latter of which had fossilized ancestors for which color was also determined. The findings were described in Nature.
Teeth
Mosasaurs possessed a thecodont dentiton, meaning that the roots were cemented deeply into the jaw bone. Mosasaurs did not use permanent teeth but instead constantly shed them. Replacement teeth developed within a pit inside the roots of the original tooth called the resorption pit. This is done through a distinctively unique eight-stage process. The first stage was characterized by the mineralization of a small tooth crown developed elsewhere that descended into the resorption pit by the second stage. In the third stage, the developing crown firmly cemented itself within the resorption pit and grew in size; by the fourth stage, it would be of the same size as the crown in the original tooth. Stages five and six were characterized by the development of the replacement tooth's root: in stage five the root developed vertically, and in stage six the root expanded in all directions to the point that the replacement tooth became exposed and actively pushed on the original tooth. In the seventh stage, the original tooth was shed and the now-independent replacement tooth began to anchor itself into the vacancy. In the eighth and final stage, the replacement tooth has grown to firmly anchor itself.
Ontogeny and growth
Mosasaur growth is not well understood, as specimens of juveniles are rare, and many were mistaken for hesperornithine birds when discovered 100 years ago. However, the discovery of several specimens of juvenile and neonate-sized mosasaurs unearthed more than a century ago indicate that mosasaurs gave birth to live young, and that they spent their early years of life out in the open ocean, not in sheltered nurseries or areas such as shallow water as previously believed. Whether mosasaurs provided parental care, like other marine reptiles such as plesiosaurs, is currently unknown. The discovery of young mosasaurs was published in the journal Palaeontology.
Possible eggs
A 2020 study published in Nature described a large fossilized hatched egg from Antarctica from the very end of the Cretaceous, about 68 million years ago. The egg is considered one of the largest amniote eggs ever known, rivalling that of the elephant bird, and due to its soft, thin, folded texture, it likely belonged to a marine animal. While the organism that produced it remains unknown, the egg's pore structure is very similar to that of extant lepidosaurs such as lizards and snakes, and presence of mosasaur fossils nearby indicates that it may have been a mosasaur egg. It is unknown whether the egg was laid on land or in the water. The egg was assigned to the newly described oospecies Antarcticoolithus bradyi. However, it has been proposed that this egg belonged to a dinosaur.
Environment
Paleontologists compared the taxonomic diversity and patterns of morphological disparity in mosasaurs with sea level, sea surface temperature, and stable carbon isotope curves for the Upper Cretaceous to explore factors that may have influenced their evolution. No single factor unambiguously accounts for all radiations, diversification, and extinctions; however, the broader patterns of taxonomic diversification and morphological disparity point to niche differentiation in a "fishing up" scenario under the influence of "bottom-up" selective pressures. The most likely driving force in mosasaur evolution was high productivity in the Late Cretaceous, driven by tectonically controlled sea levels and climatically controlled ocean stratification and nutrient delivery. When productivity collapsed at the end of the Cretaceous, coincident with bolide impact, mosasaurs became extinct.
Sea levels were high during the Cretaceous period, causing marine transgressions in many parts of the world, and a great inland seaway in what is now North America. Mosasaur fossils have been found in the Netherlands, Belgium, Denmark, Portugal, Sweden, South Africa, Spain, France, Germany, Poland, the Czech Republic, Italy Bulgaria, the United Kingdom, Russia, Ukraine, Kazakhstan, Azerbaijan, Japan, Egypt, Israel, Jordan, Syria, Turkey, Niger, Angola, Morocco, Australia, New Zealand, and on Vega Island off the coast of Antarctica. Tooth taxon Globidens timorensis is known from the island of Timor; however, the phylogenetic placement of this species is uncertain and it might not even be a mosasaur.
Mosasaurs have been found in Canada in Manitoba and Saskatchewan and in much of the contiguous United States. Complete or partial specimens have been found in Alabama, Mississippi, New Jersey, Tennessee, and Georgia, as well as in states covered by the Cretaceous seaway: Texas, southwest Arkansas, New Mexico, Kansas, Colorado, Nebraska, South Dakota, Montana, Wyoming, and the Pierre Shale/Fox Hills formations of North Dakota. Lastly, mosasaur bones and teeth are also known from Colombia, Brazil, and Chile.
Many of the so-called 'dinosaur' remains found on New Zealand are actually mosasaurs and plesiosaurs, both being Mesozoic predatory marine reptiles.
The largest mosasaur currently on public display is Bruce, a 65-70%-complete specimen of Tylosaurus pembinensis dating from the late Cretaceous Period, approximately 80 million years ago, and measuring 13.05 m (42.815 ft) from nose tip to tail tip. Bruce was discovered in 1974 north of Thornhill, Manitoba, Canada, and resides at the nearby Canadian Fossil Discovery Centre in Morden, Manitoba. Bruce was awarded the Guinness Record for the largest mosasaur on public display in 2014.
Discovery
The first publicized discovery of a partial fossil mosasaur skull in 1764 by quarry workers in a subterranean gallery of a limestone quarry in Mount Saint Peter, near the Dutch city of Maastricht, preceded any major dinosaur fossil discoveries, but remained little known. However, a second find of a partial skull drew the Age of Enlightenment's attention to the existence of fossilized animals that were different from any known living creatures. When the specimen was discovered between 1770 and 1774, Johann Leonard Hoffmann, a surgeon and fossil collector, corresponded about it with the most influential scientists of his day, making the fossil famous. The original owner, though, was Godding, a canon of Maastricht cathedral.
When the French revolutionary forces occupied Maastricht in 1794, the carefully hidden fossil was uncovered, after a reward, it is said, of 600 bottles of wine, and transported to Paris. After it had been earlier interpreted as a fish, a crocodile, and a sperm whale, the first to understand its lizard affinities was the Dutch scientist Adriaan Gilles Camper in 1799. In 1808, Georges Cuvier confirmed this conclusion, although le Grand Animal fossile de Maëstricht was not actually named Mosasaurus ('Meuse reptile') until 1822 and not given its full species name, Mosasaurus hoffmannii, until 1829. Several sets of mosasaur remains, which had been discovered earlier at Maastricht but were not identified as mosasaurs until the 19th century, have been on display in the Teylers Museum, Haarlem, procured from 1790.
The Maastricht limestone beds were rendered so famous by the mosasaur discovery, they have given their name to the final six-million-year epoch of the Cretaceous, the Maastrichtian.
Classification
Relationship with modern squamates
Lower classifications
The traditional view of mosasaur evolution held that all paddle-limbed (hydropedal) mosasaurs originated from a single common ancestor with functional legs (plesiopedal). However, this was shaken with the discovery of Dallasaurus, a plesiopedal mosasauroid more closely related to the Mosasaurinae than other mosasaurs. Bell and Polycn (2005) grouped these outside mosasaurs into two clades: the Russellosaurina, whose basal members include plesiopedal genera (Tethysaurinae) of their own and derived members consisting of the Plioplatecarpinae and Tylosaurinae; and the Halisauromorpha, containing the Halisaurinae. The placement of Dallasaurus suggested that the Russellosaurina and Halisauromorpha may have evolved a hydropedal form independently, the former through the tethysaurines, meaning that their placement within the Mosasauridae creates an unnatural polyphyly and thus potentially invalid. Caldwell informally proposed in a 2012 publication that the definition of a mosasaur must thus be redefined into one that does not consider russellosaurines and halisauromorphs as true mosasaurs, but as an independent group of marine lizards.
However, phylogenetic studies of mosasaurs can be fickle, especially when wild card taxa like Dallasaurus remain poorly understood. For example, some studies such as a 2009 analysis by Dutchak and Caldwell instead found that Dallasaurus was ancestral to both russellosaurines and mosasaurines, although results were inconsistent in later studies. A 2017 study by Simoes et al. noted that utilization of different methods of phylogenetic analyses can yield different findings and ultimately found an indication that tethysaurines were a case of hydropedal mosasaurs reversing back to a plesiopedal condition rather than an independent ancestral feature.
The following cladograms illustrate the two views of mosasaur evolution. Topology A follows an ancestral state reconstruction from an implied weighted maximum parsimony tree by Simoes et al. (2017), which contextualizes a single marine origin with tethysaurine reversal. Topologies B and C illustrate the multiple-origins hypothesis of hydropedality; the former follows Makádi et al. (2012), while the latter follows a PhD dissertation by Mekarski (2017) that experimentally includes dolichosaur and poorly-represented aigialosaur taxa. Placement of major group names follow definitions by Madzia and Cau (2017).
Phylogeny
The following diagram illustrates simplified phylogenies of the three major mosasaur groups as recovered by Strong et al. (2020), Longrich et al. (2021), and Longrich et al. (2022).
Distribution
Though no individual genus or subfamily is found worldwide, the Mosasauridae as a whole achieved global distribution during the Late Cretaceous with many locations typically having complex mosasaur faunas with multiple different genera and species in different ecological niches.
Two African countries are particularly rich in mosasaurs: Morocco and Angola.
| Biology and health sciences | Prehistoric squamates | Animals |
801330 | https://en.wikipedia.org/wiki/Gaia%20%28spacecraft%29 | Gaia (spacecraft) | Gaia is a space observatory of the European Space Agency (ESA), launched in 2013 and operated until March 2025 (planned). The spacecraft is designed for astrometry: measuring the positions, distances and motions of stars with unprecedented precision, and the positions of exoplanets by measuring attributes about the stars they orbit such as their apparent magnitude and color. The mission aims to construct by far the largest and most precise 3D space catalog ever made, totalling approximately 1 billion astronomical objects, mainly stars, but also planets, comets, asteroids and quasars, among others.
To study the precise position and motion of its target objects, the spacecraft monitored each of them about 70 times over the five years of the nominal mission (2014–2019), and about as many during its extension. Due to its detectors not degrading as fast as initially expected, the mission was given an extension. As of March 2023, the spacecraft has enough micro-propulsion fuel to operate until the second quarter of 2025. Gaia targets objects brighter than magnitude 20 in a broad photometric band that covers the extended visual range between near-UV and near infrared; such objects represent approximately 1% of the Milky Way population. Additionally, Gaia is expected to detect thousands to tens of thousands of Jupiter-sized exoplanets beyond the Solar System by using the astrometry method, 500,000 quasars outside this galaxy and tens of thousands of known and new asteroids and comets within the Solar System.
The Gaia mission continues to create a precise three-dimensional map of astronomical objects throughout the Milky Way and map their motions, which encode the origin and subsequent evolution of the Milky Way. The spectrophotometric measurements provide detailed physical properties of all stars observed, characterizing their luminosity, effective temperature, gravity and elemental composition. This massive stellar census is providing the basic observational data to analyze a wide range of important questions related to the origin, structure and evolutionary history of the Milky Way galaxy.
The successor to the Hipparcos mission (operational 1989–1993), Gaia is part of ESA's Horizon 2000+ long-term scientific program. Gaia was launched on 19 December 2013 by Arianespace using a Soyuz ST-B/Fregat-MT rocket flying from Kourou in French Guiana. The spacecraft currently operates in a Lissajous orbit around the Sun–Earth L2 Lagrangian point. The mission officially ended its operations on 15 January 2025.
History
The Gaia space telescope has its roots in ESA's Hipparcos mission (1989–1993). Its mission was proposed in October 1993 by Lennart Lindegren (Lund Observatory, Lund University, Sweden) and Michael Perryman (ESA) in response to a call for proposals for ESA's Horizon Plus long-term scientific programme. It was adopted by ESA's Science Programme Committee as cornerstone mission number 6 on 13 October 2000, and the B2 phase of the project was authorised on 9 February 2006, with EADS Astrium taking responsibility for the hardware. The name "Gaia" was originally derived as an acronym for Global Astrometric Interferometer for Astrophysics. This reflected the optical technique of interferometry that was originally planned for use on the spacecraft. While the working method evolved during studies and the acronym is no longer applicable, the name Gaia remained to provide continuity with the project.
The total cost of the mission is around €740 million (~ $1 billion), including the manufacture, launch and ground operations. Gaia was completed two years behind schedule and 16% above its initial budget, mostly due to the difficulties encountered in polishing Gaia ten silicon carbide mirrors and assembling and testing the focal plane camera system.
Objectives
The Gaia space mission has the following objectives:
To determine the intrinsic luminosity of a star requires knowledge of its distance. One of the few ways to achieve this without physical assumptions is through the star's parallax, but atmospheric effects and instrumental biases degrade the precision of parallax measurements. For instance, Cepheid variables are used as standard candles to measure distances to galaxies, but their own distances are poorly known. Thus, quantities depending on them, such as the speed of expansion of the universe, remain inaccurate.
Observations of the faintest objects will provide a more complete view of the stellar luminosity function. Gaia will observe 1 billion stars and other bodies, representing 1% of such bodies in the Milky Way galaxy. All objects up to a certain magnitude must be measured in order to have unbiased samples.
To permit a better understanding of the more rapid stages of stellar evolution (such as the classification, frequency, correlations and directly observed attributes of rare fundamental changes and of cyclical changes). This has to be achieved by detailed examination and re-examination of a great number of objects over a long period of operation. Observing a large number of objects in the galaxy is also important to understand the dynamics of this galaxy.
Measuring the astrometric and kinematic properties of a star is necessary in order to understand the various stellar populations, especially the most distant.
Spacecraft
Gaia was launched by Arianespace, using a Soyuz ST-B rocket with a Fregat-MT upper stage, from the Ensemble de Lancement Soyouz at Kourou in French Guiana on 19 December 2013 at 09:12 UTC (06:12 local time). The satellite separated from the rocket's upper stage 43 minutes after launch at 09:54 UTC.
The craft headed towards the Sun–Earth Lagrange point L2 located approximately 1.5 million kilometres from Earth, arriving there 8 January 2014. The L2 point provides the spacecraft with a very stable gravitational and thermal environment. There, it uses a Lissajous orbit that avoids blockage of the Sun by the Earth, which would limit the amount of solar energy the satellite could produce through its solar panels, as well as disturb the spacecraft's thermal equilibrium. After launch, a 10-metre-diameter sunshade was deployed. The sunshade always maintains a fixed 45 degree angle to the Sun, while precessing to scan the sky, thus keeping all telescope components cool and powering Gaia using solar panels on its surface. These factors and the materials used in its creation allow Gaia to function in conditions between -170°C and 70°C.
Scientific instruments
The Gaia payload consists of three main instruments:
The astrometry instrument (Astro) precisely determines the positions of all stars brighter than magnitude 20 by measuring their angular position. By combining the measurements of any given star over the duration of the mission, it will be possible to determine its parallax, and therefore its distance, and its proper motion—the velocity of the star projected on the plane of the sky.
The photometric instrument (BP/RP) allows the acquisition of luminosity measurements of stars over the 320–1000 nm spectral band, of all stars brighter than magnitude 20. The blue and red photometers (BP/RP) are used to determine stellar properties such as temperature, mass, age and elemental composition. Multi-colour photometry is provided by two low-resolution fused-silica prisms dispersing all the light entering the field of view in the along-scan direction prior to detection. The Blue Photometer (BP) operates in the wavelength range 330–680 nm; the Red Photometer (RP) covers the wavelength range 640–1050 nm.
The Radial-Velocity Spectrometer (RVS) is used to determine the velocity of celestial objects along the line of sight by acquiring high-resolution spectra in the spectral band 847–874 nm (field lines of calcium ion) for objects up to magnitude 17. Radial velocities are measured with a precision between 1 km/s (V=11.5) and 30 km/s (V=17.5). The measurements of radial velocities are important "to correct for perspective acceleration which is induced by the motion along the line of sight". The RVS reveals the velocity of the star along the line of sight of Gaia by measuring the Doppler shift of absorption lines in a high-resolution spectrum.
In order to maintain the fine pointing to focus on stars many light years away, the only moving parts are actuators to align the mirrors and the valves to fire the thrusters. It has no reaction wheels or gyroscopes. The spacecraft subsystems are mounted on a rigid silicon carbide frame, which provides a stable structure that will not expand or contract due to temperature. Attitude control is provided by small cold gas thrusters that can output 1.5 micrograms of nitrogen per second.
The telemetric link with the satellite is about 3 Mbit/s on average, while the total content of the focal plane represents several Gbit/s. Therefore, only a few dozen pixels around each object can be downlinked.
Measurement principles
Similar to its predecessor Hipparcos, but with a precision one hundred times greater, Gaia consists of two telescopes providing two observing directions with a fixed, wide angle of 106.5° between them. The spacecraft rotates continuously around an axis perpendicular to the two telescopes' lines of sight, with a spin period of 6 hours. Thus, every 6 hours the spacecraft scans a great circle stripe approximately 0.7 degrees wide. The spin axis in turn has a slower precession across the sky: it maintains a fixed 45 degree angle to the Sun, but follows a cone around the Sun every 63 days, giving a cycloid-like path relative to the stars. Over the course of the mission, each star is scanned many times at various scan directions, providing interlocking measurements over the full sky.
The two key telescope properties are:
1.45 × 0.5 m primary mirror for each telescope
1.0 × 0.5 m focal plane array on which light from both telescopes is projected. This in turn consists of 106 CCDs of 4500 × 1966 pixels each, for a total of 937.8 megapixels (commonly depicted as a gigapixel-class imaging device).
Each celestial object was observed on average about 70 times during the five years of the nominal mission, which has been extended to approximately ten years and will thus obtain twice as many observations. These measurements will help determine the astrometric parameters of stars: two corresponding to the angular position of a given star on the sky, two for the derivatives of the star's position over time (motion) and lastly, the star's parallax from which distance can be calculated. The radial velocity of the brighter stars is measured by an integrated spectrometer observing the Doppler effect. Because of the physical constraints imposed by the Soyuz spacecraft, Gaia focal arrays could not be equipped with optimal radiation shielding, and ESA expected their performance to suffer somewhat toward the end of the initial five-year mission. Ground tests of the CCDs while they were subjected to radiation provided reassurance that the primary mission's objectives can be met.
An atomic clock on board Gaia plays a crucial role in achieving the mission's primary objectives. Gaia rotates with angular velocity of 60"/sec or 0.6 microarcseconds in 10 nanoseconds. Therefore, in order to meet its positioning goals, Gaia must be able to record the exact time of observation to within nanoseconds. Furthermore, no systematic positioning errors over the rotational period of 6 hours should be introduced by the clock performance. For the timing error to be below 10 nanoseconds over each rotational period, the frequency stability of the on-board clock needs to be better than 10−12. The rubidium atomic clock aboard the Gaia spacecraft has a stability reaching ~ 10−13 over each rotational period of 21600 seconds.
Gaias measurements contribute to the creation and maintenance of a high-precision celestial reference frame, the Barycentric Celestial Reference System (BCRS), which is essential for both astronomy and navigation. This reference frame serves as a fundamental grid for positioning celestial objects in the sky, aiding astronomers in various research endeavors. All observations, regardless of the actual positioning of the spacecraft, must be expressed in terms of this reference system. As a fully relativistic model, the influence of the gravitational field of the solar-system must be taken into account, including such factors as the gravitational light-bending due to the Sun, the major planets and the Moon.
The expected accuracies of the final catalogue data have been calculated following in-orbit testing, taking into account the issues of stray light, degradation of the optics, and the basic angle instability. The best accuracies for parallax, position and proper motion are obtained for the brighter observed stars, apparent magnitudes 3–12. The standard deviation for these stars is expected to be 6.7 micro-arcseconds or better. For fainter stars, error levels increase, reaching 26.6 micro-arcseconds error in the parallax for 15th-magnitude stars, and several hundred micro-arcseconds for 20th-magnitude stars. For comparison, the best parallax error levels from the new Hipparcos reduction are no better than 100 micro-arcseconds, with typical levels several times larger.
Launch and orbit
In October 2013 ESA had to postpone Gaia original launch date, due to a precautionary replacement of two of Gaia transponders. These are used to generate timing signals for the downlink of science data. A problem with an identical transponder on a satellite already in orbit motivated their replacement and reverification once incorporated into Gaia. The rescheduled launch window was from 17 December 2013 to 5 January 2014, with Gaia slated for launch on 19 December.
Gaia was successfully launched on 19 December 2013 at 09:12 UTC.
About three weeks after launch, on 8 January 2014, it reached its designated orbit around the Sun-Earth L2 Lagrange point (SEL2), about 1.5 million kilometers from Earth.
In 2015, the Pan-STARRS observatory discovered an object orbiting the Earth, which the Minor Planet Center catalogued as object . It was soon found to be an accidental rediscovery of the Gaia spacecraft and the designation was promptly retracted.
Issues
Stray light problem
Shortly after launch, ESA revealed that Gaia suffers from a stray light problem. The problem was initially thought to be due to ice deposits reflecting some of the light diffracted around the edges of the sunshield into the telescope apertures and on towards the focal plane. The actual source of the stray light was later identified as the fibers of the sunshield, protruding beyond the edges of the shield. This results in a "degradation in science performance [which] will be relatively modest and mostly restricted to the faintest of Gaia one billion stars." Mitigation schemes were implemented to improve performance. The degradation is more severe for the RVS spectrograph than for the astrometry measurements, because it spreads the light of the star onto a much larger number of detector pixels, each of which collects scattered light.
This kind of problem has some historical background. In 1985 on STS-51-F, the Space Shuttle Spacelab-2 mission, another astronomical mission hampered by stray debris was the Infrared Telescope (IRT), in which a piece of mylar insulation broke loose and floated into the line-of-sight of the telescope causing corrupted data. The testing of stray-light and baffles is a noted part of space imaging instruments.
Micrometeoroid hit
In April 2024, a micrometeoroid hit and damaged Gaia's protective cover, creating "a little gap that allowed stray sunlight – around one billionth of the intensity of direct sunlight felt on Earth – to occasionally disrupt Gaia’s very sensitive sensors". In May, the electronics of one of the CCDs failed, which caused a high rate of false detections. After that, the engineers refocused Gaia'''s optics "for the final time".
Mission progress
The testing and calibration phase, which started while Gaia was en route to SEL2 point, continued until the end of July 2014, three months behind schedule due to unforeseen issues with stray light entering the detector. After the six-month commissioning period, the satellite started its nominal five-year period of scientific operations on 25 July 2014 using a special scanning mode that intensively scanned the region near the ecliptic poles; on 21 August 2014 Gaia began using its normal scanning mode which provides more uniform coverage.
Although it was originally planned to limit Gaias observations to stars fainter than magnitude 5.7, tests carried out during the commissioning phase indicated that Gaia could autonomously identify stars as bright as magnitude 3. When Gaia entered regular scientific operations in July 2014, it was configured to routinely process stars in the magnitude range 3 – 20. On the bright side of that limit, special operational procedures download raw scanning data for the remaining 230 stars brighter than magnitude 3; methods to reduce and analyse these data are being developed; and it is expected that there will be "complete sky coverage at the bright end" with standard errors of "a few dozen μas".
On 30 August 2014, Gaia discovered its first supernova in another galaxy. On 3 July 2015, a map of the Milky Way by star density was released, based on data from the spacecraft. As of August 2016, "more than 50 billion focal plane transits, 110 billion photometric observations and 9.4 billion spectroscopic observations have been successfully processed."
In 2018 the Gaia mission was extended to 2020, and in 2020 it was further extended through 2022, with an additional "indicative extension" extending through 2025. The limiting factor to further mission extensions is the supply of nitrogen for the cold gas thrusters of the micro-propulsion system. The amount of dinitrogen tetroxide (NTO) and monomethylhydrazine (MMH) for the chemical propulsion subsystem on board might be enough to stabilize the spacecraft at L2 for several decades. Without the cold gas, though, the space craft can no longer be pointed on a microarcsecond scale.
In March 2023, the Gaia mission was extended through the second quarter of 2025, when the spacecraft was expected to run out of cold gas propellant.
End of mission Gaias last targeted observation was done on 10 January 2025. After several weeks of onboard technology tests, Gaia will leave its orbit near and be put into a heliocentric orbit away from Earth's sphere of influence. After downlinking all remaining data to Earth, Gaia will be decommissioned and passivated on 27 March 2025. The mission will then enter a post-operations phase to complete and publish the final Gaia Data Release, DR5, by the end of 2030.
Data releases
Several Gaia catalogues are released over the years each time with increasing amounts of information and better astrometry; the early releases also miss some stars, especially fainter stars located in dense star fields and members of close binary pairs. The first data release, Gaia DR1, based on 14 months of observation was on 14 September 2016. The data release includes "positions and ... magnitudes for 1.1 billion stars using only Gaia data; positions, parallaxes and proper motions for more than 2 million stars" based on a combination of Gaia and Tycho-2 data for those objects in both catalogues; "light curves and characteristics for about 3,000 variable stars; and positions and magnitudes for more than 2000 ... extragalactic sources used to define the celestial reference frame".
The second data release (DR2), which occurred on 25 April 2018, is based on 22 months of observations made between 25 July 2014 and 23 May 2016. It includes positions, parallaxes and proper motions for about 1.3 billion stars and positions of an additional 300 million stars in the magnitude range g = 3–20, red and blue photometric data for about 1.1 billion stars and single colour photometry for an additional 400 million stars, and median radial velocities for about 7 million stars between magnitude 4 and 13. It also contains data for over 14,000 selected Solar System objects.
Due to uncertainties in the data pipeline, the third data release, based on 34 months of observations, has been split into two parts so that data that was ready first, was released first. The first part, EDR3 ("Early Data Release 3"), consisting of improved positions, parallaxes and proper motions, was released on 3 December 2020. The coordinates in EDR3 use a new version of the Gaia celestial reference frame (Gaia–CRF3), based on observations of 1,614,173 extragalactic sources, 2,269 of which were common to radio sources in the third revision of the International Celestial Reference Frame (ICRF3). Included is the Gaia Catalogue of Nearby Stars (GCNS), containing 331,312 stars within (nominally) .
The full DR3, published on 13 June 2022, includes the EDR3 data plus Solar System data; variability information; results for non-single stars, for quasars, and for extended objects; astrophysical parameters; and a special data set, the Gaia Andromeda Photometric Survey (GAPS).
Future releases
The full data release for the five-year nominal mission, DR4, will include full astrometric, photometric and radial-velocity catalogues, variable-star and non-single-star solutions, source classifications plus multiple astrophysical parameters for stars, unresolved binaries, galaxies and quasars, an exo-planet list and epoch and transit data for all sources. Additional release(s) will take place depending on mission extensions. Most measurements in DR4 are expected to be 1.7 times more precise than DR2; proper motions will be 4.5 times more precise. DR4 is expected to be released no earlier than mid-2026.
The final Gaia catalogue, DR5, will consist of all data collected during the lifespan of the mission. It will be 1.4 times more precise than DR4, while proper motions will be 2.8 times more precise than DR4. It will be published no earlier than the end of 2030. All data of all catalogues will be available in an online data base that is free to use.
An outreach application, Gaia Sky, has been developed to explore the galaxy in three dimensions using Gaia data.
Data processing
The overall data volume that was retrieved from the spacecraft during the nominal five-year mission at a compressed data rate of 1 Mbit/s is approximately 60 TB, amounting to about 200 TB of usable uncompressed data on the ground, stored in an InterSystems Caché database. The responsibility of the data processing, partly funded by ESA, is entrusted to a European consortium, the Data Processing and Analysis Consortium (DPAC), which was selected after its proposal to the ESA Announcement of Opportunity released in November 2006. DPAC's funding is provided by the participating countries and has been secured until the production of Gaia final catalogue.Gaia sends back data for about eight hours every day at about 5 Mbit/s. ESA's three 35-metre-diameter radio dishes of the ESTRACK network in Cebreros, Spain, Malargüe, Argentina and New Norcia, Australia, receive the data.
Significant results
In July 2017, the Gaia–ESO Survey reported using Gaia data to find double-, triple-, and quadruple- stars. Using advanced techniques they identified 342 binary candidates, 11 triple candidates, and 1 quadruple candidate. Nine of these had been identified by other means, thus confirming that the technique can correctly identify multiple star systems. The possible quadruple star system is HD 74438, which was, in a paper published in 2022, identified as a possible progenitor of a sub-Chandrasekhar Type Ia supernovae.
In November 2017, scientists led by Davide Massari of the Kapteyn Astronomical Institute, University of Groningen, Netherlands released a paper describing the characterization of proper motion (3D) within the Sculptor dwarf galaxy, and of that galaxy's trajectory through space and with respect to the Milky Way, using data from Gaia and the Hubble Space Telescope. Massari said, "With the precision achieved we can measure the yearly motion of a star on the sky which corresponds to less than the size of a pinhead on the Moon as seen from Earth." The data showed that Sculptor orbits the Milky Way in a highly elliptical orbit; it is currently near its closest approach at a distance of about , but the orbit can take it out to around distant.
In October 2018, Leiden University astronomers were able to determine the orbits of 20 hypervelocity stars from the DR2 dataset. Expecting to find a single star exiting the Milky Way, they instead found seven. More surprisingly, the team found that 13 hypervelocity stars were instead approaching the Milky Way, possibly originating from as-of-yet unknown extragalactic sources. Alternatively, they could be halo stars to this galaxy, and further spectroscopic studies will help determine which scenario is more likely. Independent measurements have demonstrated that the greatest Gaia radial velocity among the hypervelocity stars is contaminated by light from nearby bright stars in a crowded field and cast doubt on the high Gaia radial velocities of other hypervelocity stars.
In late October 2018, the galactic population Gaia-Enceladus, the remains of a major merger with the defunct Enceladus dwarf, was discovered. This system is associated with at least 13 globular clusters, and the creation of the Thick Disk of the Milky Way. It represents a significant merger about 10 billion years ago in the Milky Way Galaxy.
In November 2018, the galaxy Antlia 2 was discovered. It is similar in size to the Large Magellanic Cloud, despite being 10,000 times fainter. Antlia 2 has the lowest surface brightness of any galaxy discovered.
In December 2019 the star cluster Price-Whelan 1 was discovered. The cluster belongs to the Magellanic Clouds and is located in the leading arm of these Dwarf Galaxies. The discovery suggests that the stream of gas extending from the Magellanic Clouds to the Milky Way is about half as far from the Milky Way as previously thought.
The Radcliffe wave was discovered in data measured by Gaia, published in January 2020.
In November 2020, Gaia measured the acceleration of the solar system towards the galactic center as 0.23 nanometers/s2.
In March 2021, the European Space Agency announced that Gaia had identified a transiting exoplanet for the first time. The planet was discovered orbiting solar-type star Gaia EDR3 3026325426682637824. Following its initial discovery, the PEPSI spectrograph from the Large Binocular Telescope (LBT) in Arizona was used to confirm the discovery and categorise it as a Jovian planet, a gas planet composed of hydrogen and helium gas. In May 2022, the confirmation of this exoplanet, designated Gaia-1b, was formally published, along with a second planet, Gaia-2b.
Based on its data, Gaia's Hertzsprung–Russell diagram (HR diagram) is one of the most accurate ones ever produced of the Milky Way Galaxy.
Analysis of Gaia DR3 data in 2022 revealed a Sun-like star with the identifier Gaia DR3 4373465352415301632 orbiting a black hole, dubbed Gaia BH1. At a distance of roughly , it is the closest known black hole to Earth. Another system with a red giant orbiting a black hole, Gaia BH2, was also discovered.
In September 2023, radial velocity observations were used to confirm an exoplanet orbiting the star HIP 66074 that was first detected in Gaia DR3 astrometry data. This planet, known as HIP 66074 b or Gaia-3b, is the third Gaia exoplanet discovery to be confirmed and the first such discovery made using astrometry. In addition, another exoplanet was discovered from a gravitational microlensing event observed by Gaia, Gaia22dkv. The host star is brighter than that of any exoplanet previously detected by microlensing, potentially making the planet detectable by radial velocity as well.
In March 2024, Gaia'' discovered two streams of stars, named by researchers Shakti and Shiva, that formed more than 12 billion years ago.
GaiaNIR
GaiaNIR (Gaia Near Infra-Red) is a proposed successor of Gaia in the near-infrared. The mission would enlarge the current catalog with sources that are only (or better) visible in the near-infrared, at the cost of less precise measurements than an equivalent visible-light mission due to the broader diffraction pattern at longer wavelengths. It would at the same time improve the star parallax and particularly proper motion accuracy by revisiting the sources of the Gaia catalog. One of the main challenges in building GaiaNIR is the low technology readiness level of near-infrared time delay and integration detectors but recent progress with Avalanche photodiode detectors (APDs) is overcoming this. In a 2017 ESA report two alternative concepts using conventional near-infrared detectors and de-spin mirrors were proposed but even without the development of NIR TDI detectors the technological challenge will likely increase the cost over an ESA M-class mission and might need shared cost with other space agencies. One possible partnership with US institutions was proposed. Since then the European Space Agency Science Programme Voyage 2050 has selected the theme of "Galactic Ecosystem with Astrometry in the Near-infrared" as one of two potential L-class missions to be implemented in the coming years thus boosting the chances for GaiaNIR which proposes exactly this.
Gallery
| Technology | Space-based observatories | null |
801420 | https://en.wikipedia.org/wiki/Atmospheric%20physics | Atmospheric physics | Within the atmospheric sciences, atmospheric physics is the application of physics to the study of the atmosphere. Atmospheric physicists attempt to model Earth's atmosphere and the atmospheres of the other planets using fluid flow equations, radiation budget, and energy transfer processes in the atmosphere (as well as how these tie into boundary systems such as the oceans). In order to model weather systems, atmospheric physicists employ elements of scattering theory, wave propagation models, cloud physics, statistical mechanics and spatial statistics which are highly mathematical and related to physics. It has close links to meteorology and climatology and also covers the design and construction of instruments for studying the atmosphere and the interpretation of the data they provide, including remote sensing instruments. At the dawn of the space age and the introduction of sounding rockets, aeronomy became a subdiscipline concerning the upper layers of the atmosphere, where dissociation and ionization are important.
Remote sensing
Remote sensing is the small or large-scale acquisition of information of an object or phenomenon, by the use of either recording or real-time sensing device(s) that is not in physical or intimate contact with the object (such as by way of aircraft, spacecraft, satellite, buoy, or ship). In practice, remote sensing is the stand-off collection through the use of a variety of devices for gathering information on a given object or area which gives more information than sensors at individual sites might convey. Thus, Earth observation or weather satellite collection platforms, ocean and atmospheric observing weather buoy platforms, monitoring of a pregnancy via ultrasound, magnetic resonance imaging (MRI), positron-emission tomography (PET), and space probes are all examples of remote sensing. In modern usage, the term generally refers to the use of imaging sensor technologies including but not limited to the use of instruments aboard aircraft and spacecraft, and is distinct from other imaging-related fields such as medical imaging.
There are two kinds of remote sensing. Passive sensors detect natural radiation that is emitted or reflected by the object or surrounding area being observed. Reflected sunlight is the most common source of radiation measured by passive sensors. Examples of passive remote sensors include film photography, infrared, charge-coupled devices, and radiometers. Active collection, on the other hand, emits energy in order to scan objects and areas whereupon a sensor then detects and measures the radiation that is reflected or backscattered from the target. radar, lidar, and SODAR are examples of active remote sensing techniques used in atmospheric physics where the time delay between emission and return is measured, establishing the location, height, speed and direction of an object.
Remote sensing makes it possible to collect data on dangerous or inaccessible areas. Remote sensing applications include monitoring deforestation in areas such as the Amazon Basin, the effects of climate change on glaciers and Arctic and Antarctic regions, and depth sounding of coastal and ocean depths. Military collection during the Cold War made use of stand-off collection of data about dangerous border areas. Remote sensing also replaces costly and slow data collection on the ground, ensuring in the process that areas or objects are not disturbed.
Orbital platforms collect and transmit data from different parts of the electromagnetic spectrum, which in conjunction with larger scale aerial or ground-based sensing and analysis, provides researchers with enough information to monitor trends such as El Niño and other natural long and short term phenomena. Other uses include different areas of the earth sciences such as natural resource management, agricultural fields such as land usage and conservation, and national security and overhead, ground-based and stand-off collection on border areas.
Radiation
Atmospheric physicists typically divide radiation into solar radiation (emitted by the sun) and terrestrial radiation (emitted by Earth's surface and atmosphere).
Solar radiation contains variety of wavelengths. Visible light has wavelengths between 0.4 and 0.7 micrometers. Shorter wavelengths are known as the ultraviolet (UV) part of the spectrum, while longer wavelengths are grouped into the infrared portion of the spectrum. Ozone is most effective in absorbing radiation around 0.25 micrometers, where UV-c rays lie in the spectrum. This increases the temperature of the nearby stratosphere. Snow reflects 88% of UV rays, while sand reflects 12%, and water reflects only 4% of incoming UV radiation. The more glancing the angle is between the atmosphere and the sun's rays, the more likely that energy will be reflected or absorbed by the atmosphere.
Terrestrial radiation is emitted at much longer wavelengths than solar radiation. This is because Earth is much colder than the sun. Radiation is emitted by Earth across a range of wavelengths, as formalized in Planck's law. The wavelength of maximum energy is around 10 micrometers.
Cloud physics
Cloud physics is the study of the physical processes that lead to the formation, growth and precipitation of clouds. Clouds are composed of microscopic droplets of water (warm clouds), tiny crystals of ice, or both (mixed phase clouds). Under suitable conditions, the droplets combine to form precipitation, where they may fall to the earth. The precise mechanics of how a cloud forms and grows is not completely understood, but scientists have developed theories explaining the structure of clouds by studying the microphysics of individual droplets. Advances in radar and satellite technology have also allowed the precise study of clouds on a large scale.
Atmospheric electricity
Atmospheric electricity is the term given to the electrostatics and electrodynamics of the atmosphere (or, more broadly, the atmosphere of any planet). The Earth's surface, the ionosphere, and the atmosphere is known as the global atmospheric electrical circuit. Lightning discharges 30,000 amperes, at up to 100 million volts, and emits light, radio waves, X-rays and even gamma rays. Plasma temperatures in lightning can approach 28,000 kelvins and electron densities may exceed 1024/m3.
Atmospheric tide
The largest-amplitude atmospheric tides are mostly generated in the troposphere and stratosphere when the atmosphere is periodically heated as water vapour and ozone absorb solar radiation during the day. The tides generated are then able to propagate away from these source regions and ascend into the mesosphere and thermosphere. Atmospheric tides can be measured as regular fluctuations in wind, temperature, density and pressure. Although atmospheric tides share much in common with ocean tides they have two key distinguishing features:
i) Atmospheric tides are primarily excited by the Sun's heating of the atmosphere whereas ocean tides are primarily excited by the Moon's gravitational field. This means that most atmospheric tides have periods of oscillation related to the 24-hour length of the solar day whereas ocean tides have longer periods of oscillation related to the lunar day (time between successive lunar transits) of about 24 hours 51 minutes.
ii) Atmospheric tides propagate in an atmosphere where density varies significantly with height. A consequence of this is that their amplitudes naturally increase exponentially as the tide ascends into progressively more rarefied regions of the atmosphere (for an explanation of this phenomenon, see below). In contrast, the density of the oceans varies only slightly with depth and so there the tides do not necessarily vary in amplitude with depth.
Note that although solar heating is responsible for the largest-amplitude atmospheric tides, the gravitational fields of the Sun and Moon also raise tides in the atmosphere, with the lunar gravitational atmospheric tidal effect being significantly greater than its solar counterpart.
At ground level, atmospheric tides can be detected as regular but small oscillations in surface pressure with periods of 24 and 12 hours. Daily pressure maxima occur at 10 a.m. and 10 p.m. local time, while minima occur at 4 a.m. and 4 p.m. local time. The absolute maximum occurs at 10 a.m. while the absolute minimum occurs at 4 p.m. However, at greater heights the amplitudes of the tides can become very large. In the mesosphere (heights of ~ 50 – 100 km) atmospheric tides can reach amplitudes of more than 50 m/s and are often the most significant part of the motion of the atmosphere.
Aeronomy
Aeronomy is the science of the upper region of the atmosphere, where dissociation and ionization are important. The term aeronomy was introduced by Sydney Chapman in 1960. Today, the term also includes the science of the corresponding regions of the atmospheres of other planets. Research in aeronomy requires access to balloons, satellites, and sounding rockets which provide valuable data about this region of the atmosphere. Atmospheric tides play an important role in interacting with both the lower and upper atmosphere. Amongst the phenomena studied are upper-atmospheric lightning discharges, such as luminous events called red sprites, sprite halos, blue jets, and elves.
Centers of research
In the UK, atmospheric studies are underpinned by the Met Office, the Natural Environment Research Council and the Science and Technology Facilities Council. Divisions of the U.S. National Oceanic and Atmospheric Administration (NOAA) oversee research projects and weather modeling involving atmospheric physics. The US National Astronomy and Ionosphere Center also carries out studies of the high atmosphere. In Belgium, the Belgian Institute for Space Aeronomy studies the atmosphere and outer space. In France, there are several public or private entities researching the atmosphere, as an example météo-France (Météo-France), several laboratories in the national scientific research center (such as the laboratories in the IPSL group).
| Physical sciences | Physics basics: General | Physics |
802029 | https://en.wikipedia.org/wiki/Warp%20and%20weft | Warp and weft | In the manufacture of cloth, warp and weft are the two basic components in weaving to transform thread and yarn into textile fabrics. The vertical warp yarns are held stationary in tension on a loom (frame) while the horizontal weft (also called the woof) is drawn through (inserted over and under) the warp thread. In the terminology of weaving, each warp thread is called a warp end (synonymous terms are fill yarn and filling yarn); a pick is a single weft thread that crosses the warp thread.
In the 18th century, the Industrial Revolution facilitated the industrialisation of the production of textile fabrics with the "picking stick" and the "flying shuttle", which latter was invented by John Kay, in 1733. The mechanised power loom was patented by Edmund Cartwright in 1785, which allowed sixty picks per minute.
Etymology
The word weft derives from the Old English word , to weave. Warp means "that across which the woof is thrown". (Old English , from , to throw, cf. German , Dutch ).
Warp
The warp is the set of yarns or other things stretched in place on a loom before the weft is introduced during the weaving process. It is regarded as the longitudinal set in a finished fabric with two or more sets of elements.
The term is also used for a set of yarns established before the interworking of weft yarns by some other method, such as finger manipulation, yielding wrapped or twined structures. Very simple looms use a spiral warp, in which the warp is made up of a single, very long yarn wound in a spiral pattern around a pair of sticks or beams.
The warp must be strong to be held under high tension during the weaving process, unlike the weft which carries almost no tension. This requires the yarn used for warp ends, or individual warp threads, to be made of spun and plied fibre. Traditionally natural fibres such as wool, linen, alpaca, and silk were used. However, improvements in spinning technology during the Industrial Revolution created cotton yarn of sufficient strength to be used in mechanized weaving. Later, synthetic fibres such as nylon or rayon were employed.
While most weaving is weft-faced, warp-faced textiles are created using densely arranged warp threads. In these the design is in the warp, requiring all colors to be decided upon and placed during the first part of the weaving process, which cannot be changed. Such limitations of color placement create weavings defined by length-wise stripes and vertical designs. Many South American cultures, including the ancient Incas and Aymaras, employed backstrap weaving, which uses the weight of the weaver's body to control the tension of the loom.
Weft
Because the weft does not have to be stretched on a loom the way the warp is, it can generally be less strong. It is usually made of spun fibre, originally wool, flax and cotton, today often of synthetic fibre such as nylon or rayon.
The weft is threaded through the warp using a "shuttle", air jets or "rapier grippers". Handlooms were the original weaver's tool, with the shuttle being threaded through alternately raised warps by hand.
As metaphor
The expression "warp and weft" (also "warp and woof" and "woof and warp") is used metaphorically the way "fabric" is; e.g., "the warp and woof of a student's life" equates to "the fabric of a student's life". Warp and weft are sometimes used even more generally in literature to describe the basic dichotomy of the world we live in, as in, up/down, in/out, black/white, Sun/Moon, yin/yang, etc. The expression is also used similarly for the underlying structure upon which something is built. The terms "warp" and "woof" are also found in some English translations of the Bible in the discussion of mildews found in cloth materials in Leviticus 13:48-59.
In computing, a warp is a term for a block of parallel threads executed on a GPU or similar SIMD device.
| Technology | Weaving | null |
802146 | https://en.wikipedia.org/wiki/Rhinoceros | Rhinoceros | A rhinoceros ( ; ; ; : rhinoceros or rhinoceroses), commonly abbreviated to rhino, is a member of any of the five extant species (or numerous extinct species) of odd-toed ungulates (perissodactyls) in the family Rhinocerotidae; it can also refer to a member of any of the extinct species of the superfamily Rhinocerotoidea. Two of the extant species are native to Africa, and three to South and Southeast Asia.
Rhinoceroses are some of the largest remaining megafauna: all weigh over half a tonne in adulthood. They have a herbivorous diet, small brains for mammals of their size, one or two horns, and a thick , protective skin formed from layers of collagen positioned in a lattice structure. They generally eat leafy material, although their ability to ferment food in their hindgut allows them to subsist on more fibrous plant matter when necessary. Unlike other perissodactyls, the two African species of rhinoceros lack teeth at the front of their mouths; they rely instead on their lips to pluck food.
Rhinoceroses are killed by poachers for their horns, which are bought and sold on the black market for high prices, leading to most living rhinoceros species being considered endangered. The contemporary market for rhino horn is overwhelmingly driven by China and Vietnam, where it is bought by wealthy consumers to use in traditional Chinese medicine, among other uses. Rhino horns are made of keratin, the same material as hair and fingernails, and there is no good evidence of any health benefits. A market also exists for rhino horn dagger handles in Yemen, which was the major source of demand for rhino horn in the 1970s and 1980s.
Taxonomy and naming
The word rhinoceros is derived through Latin from the , which is composed of (rhino-, "of the nose") and (, "horn") with a horn on the nose. The name has been in use since the 14th century.
The family Rhinocerotidae consists of only four extant genera: Ceratotherium (white rhinoceros), Diceros (black rhinoceros), Dicerorhinus (Sumatran rhinoceros), and Rhinoceros (Indian and Javan rhinoceros). The living species fall into three categories. The two African species, the white rhinoceros and the black rhinoceros, belong to the tribe Dicerotini, which originated in the middle Miocene, about 14.2 million years ago. The species diverged during the early Pliocene (about 5 million years ago). The main difference between black and white rhinos is the shape of their mouths – white rhinos have broad flat lips for grazing, whereas black rhinos have long pointed lips for eating foliage. There are two living Rhinocerotini species, the Indian rhinoceros and the Javan rhinoceros, which diverged from one another about 10 million years ago. The Sumatran rhinoceros is the only surviving representative of the Dicerorhinini.
A subspecific hybrid white rhino (Ceratotherium s. simum × C. s. cottoni) was bred at the Dvůr Králové Zoo (Zoological Garden Dvur Kralove nad Labem) in the Czech Republic in 1977. Interspecific hybridisation of black and white rhinoceroses has also been confirmed.
While the black rhinoceros has 84 chromosomes (diploid number, 2N, per cell), all other rhinoceros species have 82 chromosomes. Chromosomal polymorphism might lead to varying chromosome counts. For instance, in a study there were three northern white rhinoceroses with 81 chromosomes.
Anatomy
Rhinoceroses are among the largest living land animals, with living species ranging in average weight from in the Sumatran rhinoceros, to in the white rhinoceros. Some extinct rhinocerotids were considerably smaller and larger than living rhinoceroses, with the genus Menoceras from the Early Miocene of North America having an estimated body mass of , comparable to sheep, or a pig, while Elasmotherium sibiricum from the Pleistocene of Eurasia has an estimated body mass of approximately . The skulls of rhinoceroses are generally saddle-shaped and low, with rhinoceroses being primitively characterised by the presence of a chisel-shaped upper first incisor (I1) and a tusk-like lower second incisor (i2), with all other incisors and the canines typically being lost. Black and white rhinoceroses completely lack incisors. Living rhinoceroses have either one or two horns, which are formed from columns of densely packed corneocytes originating from dermal papillae. The development and growth of rhinoceros horns is similar to that of human nails, with both being largely made of keratin. The horns are attached to a rugose (roughly textured) area on the surface of the skull. Horns are not a universal feature of rhinocerotids, with horns thought to be absent in many extinct rhinocerotids (such as most members of the subfamily Aceratheriinae). The brains of rhinoceroses are relatively small compared to body size, around in an adult black rhinoceros. The limb bones tend to be robust (proportionally thick and stocky). All living and the vast majority of extinct rhinoceroses have three toes on each foot. The body is covered in an armour of thick skin made of a dense crosslinked network of collagen fibres that is stronger and stiffer than those of other mammals. The skin exhibits prominent folding. The skin in living species is grey to brown in colour, and typically sparsely covered in hair or hairless as adults, with the exception of the eyelashes, ears, and the tail-brush. The exception is the Sumatran rhinoceros, which is often covered with a considerable amount of hair.
Behaviour and ecology
Living rhinoceroses gregariousness varies between species. Adult males tend to be solitary, and this is also true of female Asian rhinoceroses, though the females of African species sometimes form groups, with these groups being more common in white than black rhinoceroses. Rhinoceroses have widely varying diets ranging from strict grazing (such as the white rhinoceros), to largely browsing (such as the black rhinoceros) to a mixture between both (the Sumatran and Javan rhinoceros). As bulk feeders of low quality vegetation, rhinoceroses spend a majority of their time foraging. Rhinoceroses are hindgut fermenters.
All living rhinoceroses have a polyandrous and polygnous mating system where both males and females seek to mate with multiple individuals of the opposite sex. Male rhinoceroses guard reproductive age females until they are in full estrous though the females sometimes may drive away males until they are receptive. Male rhinoceroses taste the urine of female rhinoceroses and perform a flehmen response with the upper lip to determine their reproductive status. Adult males in the vicinity of oestrous females may become aggressive towards other males. These confrontations can range from ritualized behaviour, to serious fighting that can result in significant injuries. In some species, male rhinoceroses are territorial, while in other species they are not, or are only territorial depending on local environmental conditions. Females will sometimes reject males they consider undesirable, which results in them fleeing or fighting the male if cornered. During copulation, the male slides his neck up the back of the female, before using his neck as a lever to get his forelegs off the ground, before moving the front legs behind the shoulders of the female. Copulation can last several hours. Pregnancy lasts for over a year, around 460 days in the black rhinoceros and 504 days in the white rhinoceros. The female generally gives birth in a secluded area, and becomes aggressive towards other rhinoceroses for a while after giving birth. Calves typically stand up within 30 minutes of birth and begin to suck on their mother's teats within two hours of birth. The mother generally has a strong bond with her most recently born calf. The calf generally remains close to its mother the majority of the time, although at least in some species they are sometimes left considerable distances away. Up until they are around three years old, juvenile rhinoceroses are vulnerable to predation. Mothers are vigorously protective of their calves against potential predators. Juvenile one-horned rhinoceroses are rejected by their mothers around the time of the birth of her next calf. There is generally a gap of several years between females giving birth again after having her previous calf, though the gap can be as short as a year and a half. Rhinoceroses become sexually mature at around five to eight years of age, generally around a year later in males than in females in black and Sumatran rhinoceroses, though male white rhinoceroses become socio-sexually mature at around 12 years of age, four years after females start giving birth.
Extant species
White
There are two subspecies of white rhinoceros: the southern white rhinoceros (Ceratotherium simum simum) and the northern white rhinoceros (Ceratotherium simum cottoni). As of 2013, the southern subspecies has a wild population of 20,405—making them the most abundant rhino subspecies in the world. The northern subspecies is critically endangered, with all that is known to remain being two captive females. There is no conclusive explanation of the name "white rhinoceros". A popular idea that "white" is a distortion of either the Afrikaans word or the Dutch word (or its other possible spellings , , etc.,), meaning "wide" and referring to the rhino's square lips, is not supported by linguistic studies.
The white rhino has an immense body and large head, a short neck and broad chest. Females weigh and males on average, though exceptional specimens can reportedly weigh up to . The head-and-body length is and the shoulder height is . On its snout it has two horns. The front horn is larger than the other horn and averages in length and can reach . The white rhinoceros also has a prominent muscular hump that supports its relatively large head. The colour of this animal can range from yellowish brown to slate grey. Most of its body hair is found on the ear fringes and tail bristles, with the rest distributed rather sparsely over the rest of the body. White rhinos have the distinctive flat broad mouth that is used for grazing.
Black
The name "black rhinoceros" (Diceros bicornis) was chosen to distinguish this species from the white rhinoceros (Ceratotherium simum). This can be confusing, as the two species are not truly distinguishable by color. There are four subspecies of black rhino: South-central (Diceros bicornis minor), the most numerous, which once ranged from central Tanzania south through Zambia, Zimbabwe and Mozambique to northern and eastern South Africa; South-western (Diceros bicornis occidentalis) which are better adapted to the arid and semi-arid savannas of Namibia, southern Angola, western Botswana and western South Africa; East African (Diceros bicornis michaeli), primarily in Tanzania; and West African (Diceros bicornis longipes) which was declared extinct in November 2011. The native Tswanan name keitloa describes a South African variation of the black rhino in which the posterior horn is equal to or longer than the anterior horn.
An adult black rhinoceros stands high at the shoulder and is in length. An adult weighs from , exceptionally to , with the females being smaller than the males. Two horns on the skull are made of keratin with the larger front horn typically long, exceptionally up to . Sometimes, a third smaller horn may develop. The black rhino is much smaller than the white rhino, and has a pointed mouth, which it uses to grasp leaves and twigs when feeding.
During the latter half of the 20th century, their numbers were severely reduced from an estimated 70,000 in the late 1960s to a record low of 2,410 in 1995. Since then, numbers have been steadily increasing at a continental level with numbers doubling to 4,880 by the end of 2010. As of 2008, the numbers are still 90% lower than three generations ago.
Indian
The Indian rhinoceros, or greater one-horned rhinoceros, (Rhinoceros unicornis) has a single horn 20 to 60 cm long. It is nearly as large as the African white rhino. Its thick, silver-brown skin folds into the shoulder, back, and rump, giving it an armored appearance. Its upper legs and shoulders are covered in wart-like bumps, and it has very little body hair. Grown males are larger than females in the wild, weighing from . Shoulder height is . Females weigh about and are long. The record-sized specimen was approximately .
Indian rhinos once inhabited many areas ranging from Pakistan to Myanmar and maybe even parts of China. Because of humans, they now exist in only several protected areas of India (in Assam, West Bengal, and a few pairs in Uttar Pradesh) and Nepal, plus a pair in Lal Suhanra National Park in Pakistan reintroduced there from Nepal. They are confined to the tall grasslands and forests in the foothills of the Himalayas. Two-thirds of the world's Indian rhinoceroses are now confined to the Kaziranga National Park situated in the Golaghat district of Assam, India.
Javan
The Javan rhinoceros (Rhinoceros sondaicus) is one of the most endangered large mammals in the world. According to 2015 estimates, only about 60 remain, in Java, Indonesia, all in the wild. It is also the least known rhino species. Like the closely related, and larger, Indian rhinoceros, the Javan rhino has a single horn. Its hairless, hazy gray skin falls into folds into the shoulder, back, and rump, giving it an armored appearance. Its length reaches including the head, and its height . Adults are variously reported to weigh or . Male horns can reach in length, while in females they are knobs or altogether absent. These animals prefer dense lowland rain forest, tall grass and reed beds that are plentiful with large floodplains and mud wallows.
Though once widespread throughout Asia, by the 1930s, they were nearly hunted to extinction in Nepal, India, Burma, Peninsular Malaysia, and Sumatra for the supposed medical powers of their horns and blood. As of 2015, only 58–61 individuals remain in Ujung Kulon National Park, Java, Indonesia. The last known Javan rhino in Vietnam was reportedly killed for its horn in 2011 by Vietnamese poachers. Now only Java contains the last Javan rhinos.
Sumatran
The Sumatran rhinoceros (Dicerorhinus sumatrensis) is the smallest extant rhinoceros species, as well as the one with the most hair. It can be found at very high altitudes in Borneo and Sumatra. Because of habitat loss and poaching, their numbers have declined, and it has become the second most threatened rhinoceros. About 275 Sumatran rhinos are believed to remain. There are three subspecies of Sumatran rhinoceros: the Sumatran rhinoceros proper (Dicerorhinus sumatrensis sumatrensis), the Bornean rhinoceros (Dicerorhinus sumatrensis harrissoni) and the possibly extinct Northern Sumatran rhinoceros (Dicerorhinus sumatrensis lasiotis).
A mature rhino typically stands about high at the shoulder, has a length of and weighs around , though the largest individuals have been known to weigh as much as . Like the African species, it has two horns; the larger is the front (), with the smaller usually less than long. Males have much larger horns than the females. Hair can range from dense (the densest hair in young calves) to sparse. The color of these rhinos is reddish brown. The body is short and has stubby legs. The lip is prehensile.
Sumatran rhinoceros once were spread across South-east Asia, but now are on the verge of extinction, confined to several parts of Indonesia and Malaysia by reproductive isolation. There were 320 D. sumatrensis in 1995, which, by 2011, had dwindled to 216. It has been found through DNA comparison that the Sumatran rhinoceros is the most ancient extant rhinoceros and related to the extinct Eurasian woolly rhino species, Coelodonta. In 1994, Alan Rabinowitz publicly denounced governments, non-governmental organizations, and other institutions for lacking in their attempts to conserve the Sumatran rhinoceros. To conserve it, they would have to relocate them from small forests to breeding programs that could monitor their breeding success. To boost reproduction, the Malaysian and Indonesian governments could also agree to exchange the gametes of the Sumatran and (smaller) Bornean subspecies. The Indonesian and Malaysian governments have also proposed a single management unit for these two ancient subspecies.
Plantations for palm oil have taken out the living areas and led to the eradication of the rhino in Sumatra.
Evolution
Rhinocerotoidea diverged from other perissodactyls (odd-toed ungulates) by the early Eocene around 51-58 million years ago, with the earliest representatives appearing during the early-middle Eocene in Asia, around 54 million years ago. The family of modern rhinoceroses, Rhinocerotidae appeared during the middle-late Eocene around 39-40 million years ago, roughly at the same in North America and Asia, with rhinoceroses migrating into Europe at the Eocene-Oligocene boundary ~34 million years ago as part of the "Grande Coupure" along with many other Asian migrants. Rhinocerotids represented the only living family of rhinocerotoids following the end of the Oligocene epoch around 23 million and the extinction of other rhinocertoid groups such as the giant paraceratheres. During the early Miocene epoch, around 20 million years ago rhinocerotids migrated into Africa following its connection to Eurasia. The last common ancestor of living rhinoceroses (which belong to the subgroup Rhinocerotina) is thought to have lived during the Miocene, at least 15-16 million years ago. Rhinocerotids reached maximum diversity in Europe and Asia during the Miocene epoch, with the group declining in diversity during the late Miocene following unfavourable climatic change, becoming entirely extinct in North America at the beginning of the Pliocene, around 5 million years ago.
The earliest remains of the genus Rhinoceros (which includes the living Indian/one horned and Javan rhinoceros) are known from the Late Miocene, represented by remains such as an indeterminate species found in deposits in Myanmar dating to around 8-9 million years ago, with the two modern species appearing during the Early-Middle Pleistocene epoch. The earliest unambiguous relatives of white and black rhinoceros belonging to the genera Ceratotherium and Diceros, first appear during the late Miocene, with the first unambiguous appearance of modern white and black rhinoceros during the Early Pleistocene. The earliest unambiguous remains of Dicerorhinus are known from the latest Pliocene, with the appearance of the modern Sumatran rhinoceros during the Early Pleistocene.
Alongside the extant species, four additional species of rhinoceros survived into the Last Glacial Period: the woolly rhinoceros (Coelodonta antiquitatis), Elasmotherium sibiricum and two species of Stephanorhinus, Merck's rhinoceros (Stephanorhinus kirchbergensis) and the Narrow-nosed rhinoceros (Stephanorhinus hemitoechus).
Cladogram showing the relationships of recent and Late Pleistocene rhinoceros species (minus Stephanorhinus hemitoechus) based on whole nuclear genomes, after Liu et al., 2021:
denotes extinct taxa
Family Rhinocerotidae
†Teletaceras
†Uintaceras
†Epiaceratherium
†Trigonias 37–34 Ma
†Ronzotherium 37–23 Ma
†Diceratherium 33.9–11.6 Ma
†Menoceras 23.03–16.3 Ma
Subfamily Rhinocerotinae
Tribe Aceratheriini
†Aceratherium lived from 33.9 to 3.4 Ma
†Acerorhinus 13.6–7.0 Ma
†Alicornops 13.7–5.3 Ma
†Aphelops 20.43–5.33 Ma
†Chilotheridium 23.0–11.6 Ma
†Chilotherium 13.7–3.4 Ma
†Floridaceras 20.4–16.3 Ma
†Hoploaceratherium 16.9–16.0 Ma
†Mesaceratherium
†Peraceras 20.6–10.3 Ma
†Plesiaceratherium 20.0–11.6 Ma
†Shansirhinus
†Sinorhinus
†Subchilotherium
Tribe Teleoceratini
†Aprotodon 28.4–5.330 Ma
†Brachydiceratherium
†Brachypotherium 20.0–5.33 Ma
†Diaceratherium 28.4–16.0 Ma
†Prosantorhinus 16.9–7.25 Ma
†Shennongtherium
†Teleoceras 16.9–4.9 Ma
Rhinocerotina Burdigalian–Present
Tribe Rhinocerotini 40.4–11.1 Ma–Present
†Gaindatherium 11.6–11.1 Ma
Subtribe Rhinocerotina 17.5 Ma–Present
†Nesorhinus .70 Ma
†Rusingaceros 17.5 Ma
Rhinoceros – Indian & Javan rhinoceros
Tribe Dicerorhinini
†Pliorhinus 5–2.5 Ma
†Coelodonta – Woolly rhinoceros
Dicerorhinus – Sumatran rhinoceros
†Dihoplus 11.610–1.810 Ma
†Stephanorhinus 9.7–0.04 Ma – Merck's rhinoceros & Narrow-nosed rhinoceros
Tribe Dicerotini 23.0–Present
Ceratotherium – White rhinoceros 7.25–Present
Diceros – Black rhinoceros 23.0–Present
†Paradiceros 15.97–11.61 Ma
†Miodiceros 11.6–5 Ma
Rhinocerotinae incertae sedis
†Protaceratherium
†Lartetotherium 15.97–8.7 Ma
Subfamily Elasmotheriinae
†Gulfoceras? 23.03–20.43 Ma
†Victoriaceros? 15 Ma
†Penetrigonias?
†Subhyracodon? 38.0–26.3 Ma
Tribe Elasmotheriini 20.0–0.1 Ma
†Bugtirhinus 20.0–16.9 Ma
†Caementodon
†Elasmotherium – Giant rhinoceros 3.6–0.039 Ma
†Hispanotherium synonymized with Huaqingtherium 16.0–7.25 Ma
†Iranotherium
†Kenyatherium
†Meninatherium
†Ningxiatherium
†Ougandatherium 20.0–16.9 Ma
†Parelasmotherium
†Procoelodonta
†Sinotherium 9.0–5.3 Ma
Predators, poaching and hunting
Adult rhinoceroses have no real predators in the wild, other than humans. Young rhinos sometimes fall prey to big cats, crocodiles, African wild dogs, and hyenas.
Although rhinos are large and aggressive and have a reputation for being resilient, they are very easily poached; they visit water holes daily and can be easily killed while they drink. As of December 2009, poaching increased globally while efforts to protect the rhino are considered increasingly ineffective. The most serious estimate, that only 3% of poachers are successfully countered, is reported of Zimbabwe, while Nepal has largely avoided the crisis. Poachers have become more sophisticated. South African officials have called for urgent action against poaching after poachers killed the last female rhino in the Krugersdorp Game Reserve near Johannesburg. Statistics from South African National Parks show that 333 rhinoceroses were killed in South Africa in 2010, increasing to 668 by 2012, over 1,004 in 2013, and over 1,338 killed in 2015. In some cases rhinos are tranquilized and their horns removed leaving them to bleed to death, while in other instances more than the horn is taken.
The Namibian government has supported the practice of rhino trophy hunting as a way to raise money for conservation. Hunting licenses for five Namibian black rhinos are auctioned annually, with the money going to the government's Game Products Trust Fund. Some conservationists and members of the public oppose or question this practice.
Horn use
Rhinoceros horns develop from subcutaneous tissues, and are made of keratinous mineralized compartments. The horns root in a germinative layer.
Rhinoceros horns are used in traditional medicines in parts of Asia, and for dagger handles in Yemen and Oman. Esmond Bradley Martin has reported on the trade for dagger handles in Yemen, which was historically a major source for the demand for rhino horn in the late 20th century. In Europe, it was historically believed that rhino horns could purify water and could detect poisoned liquids, and likely believed to be an aphrodisiac and an antidote to poison.
It is a common misconception that rhinoceros horn in powdered form is used as an aphrodisiac or a cure for cancer in traditional Chinese medicine (TCM) as Cornu Rhinoceri Asiatici (犀角, xījiǎo, "rhinoceros horn"); no TCM text in history has ever mentioned such prescriptions. In TCM, rhino horn is sometimes prescribed for fevers and convulsions, a treatment not supported by evidence-based medicine: this treatment has been compared to consuming fingernail clippings in water. In a 2021 survey of Chinese users of rhinoceros horn TCM products, the vast majority of respondents cited "dispelling heat" and "detoxification" as reasons for using rhino horn. In 1993, China signed the CITES treaty and removed rhinoceros horn from the Chinese medicine pharmacopeia, administered by the Ministry of Health. In 2011, the Register of Chinese Herbal Medicine in the United Kingdom issued a formal statement condemning the use of rhinoceros horn. A growing number of TCM educators are also speaking out against the practice, although some TCM practitioners still believe that it is a life-saving medicine.
Vietnam reportedly has the biggest number of rhino horn consumers, with their demand driving most of the poaching, which has risen to record levels. The "Vietnam CITES Management Authority" has claimed that Hanoi recently experienced a 77% drop in the usage of rhino horn, but National Geographic has challenged these claims, noticing that there was no rise in the numbers of criminals who were apprehended or prosecuted. South African rhino poaching's main destination market is Vietnam. An average sized horn can bring in as much as a quarter of a million dollars in Vietnam and many rhino range states have stockpiles of rhino horn.
Horn trade
International trade in rhinoceros horn has been declared illegal by the Convention on International Trade in Endangered Species of Wild Fauna and Flora (CITES) since 1977. A proposal by Swaziland to lift the international ban was rejected in October 2016. Domestic sale of rhinoceros horn in South Africa, home of 80% of the remaining rhino population, was banned as of 2009. The ban was overturned in a court case in 2017, and South Africa plans to draft regulations for the sale of rhino horn, possibly including export for "non-commercial purposes". The South African government has proposed that a legal trade of rhino horn be established, arguing that this could reduce poaching and prevent the extinction of this species.
In March 2013, some researchers suggested that the only way to reduce poaching would be to establish a regulated trade based on humane and renewable harvesting from live rhinos. The World Wildlife Fund opposes legalization of the horn trade, as it may increase demand, while IFAW released a report by EcoLarge, suggesting that more thorough knowledge of economic factors is required to justify the pro-trade option.
Conservation
According to the World Wide Fund for Nature, conservation of African rhinoceroses as consumers of large amounts of vegetation is crucial to maintaining the shape of the African landscape and the natural resources of local communities.
Ways to prevent poaching
Horn removal
To prevent poaching, in certain areas, rhinos have been tranquillized and their horns removed. Armed park rangers, particularly in South Africa, are also working on the front lines to combat poaching, sometimes killing poachers who are caught in the act. A 2012 spike in rhino killings increased concerns about the future of the species.
Horn poisoning
In 2011, the Rhino Rescue Project began a horn-trade control method consisting of infusing the horns of living rhinos with a mixture of a pink dye and an acaricide (to kill ticks) which is safe for rhinos but toxic to humans. The procedure also includes inserting three RFID identification chips and taking DNA samples. Because of the fibrous nature of rhino horn, the pressurized dye infuses the interior of the horn but does not color the surface or affect rhino behavior. Depending on the quantity of horn a person consumes, experts believe the acaricide would cause nausea, stomach-ache, and diarrhea, and possibly convulsions. It would not be fatal—the primary deterrent is the knowledge that the treatment has been applied, communicated by signs posted at the refuges. The original idea grew out of research into the horn as a reservoir for one-time tick treatments, and experts selected an acaricide they think is safe for the rhino, oxpeckers, vultures, and other animals in the preserve's ecosystem. Proponents claim that the dye cannot be removed from the horns, and remains visible on X-ray scanners even when the horn is ground to a fine powder.
The UK charity organization Save the Rhino has criticized horn poisoning on moral and practical grounds. The organization questions the assumptions that the infusion technique works as intended, and that even if the poison were effective, whether middlemen in a lucrative, illegal trade would care much about the effect it would have on buyers. Additionally, rhino horn is increasingly purchased for decorative use, rather than for use in traditional medicine. Save the Rhino questions the feasibility of applying the technique to all African rhinos, since workers would have to reapply the acaricide every four years. It was also reported that one out of 150 rhinos treated did not survive the anesthesia.
Artificial substitute for rhinoceros horn
Another way to undercut the rhinoceros horn market has been suggested by Matthew Markus of Pembient, a biotechnology firm. He proposes the synthesis of an artificial substitute for rhinoceros horn. To enable authorities to distinguish the bioengineered horn from real rhinoceros horn, the genetic code of the bioengineered horn could be registered, similar to the DNA of living rhinoceros in the RhODIS (Rhino DNA Index System). Initial responses from many conservationists were negative, but a 2016 report from TRAFFIC—which monitors trade in wildlife and animal parts—conceded that it "...would be rash to rule out the possibility that trade in synthetic rhinoceros horn could play a role in future conservation strategies".
Historical representations
Woolly rhinoceroses are depicted the European Paleolithic art, such as in cave paintings in Chauvet Cave in France, which date to around 30-40,000 years ago.
Greek historian and geographer Agatharchides (2nd century BC) mentions the rhinoceros in his book On the Erythraean Sea.
In Khmer art, the Hindu god Agni is depicted with a rhinoceros as his vahana. Similarly in medieval era Thai literature, Agni also called Phra Phloeng is sometimes described as riding a rhinoceros.
Albrecht Dürer created a famous woodcut of a rhinoceros in 1515, based on a written description and brief sketch by Valentim Fernandes, a German printer resident in Lisbon. He never saw the animal itself, so Dürer's Rhinoceros is a somewhat inaccurate depiction.
There are legends about rhinoceroses stamping out fire in Burma, India, and Malaysia. The mythical rhinoceros has a special name in Malay, , wherein means rhinoceros, and means fire. The animal would come when a fire was lit in the forest and stamp it out. There are no recent confirmations of this phenomenon. This legend was depicted in the film The Gods Must Be Crazy (1980), which shows an African rhinoceros putting out two campfires.
In 1974, a lavender rhinoceros symbol began to be used as a symbol of the gay community in Boston, United States.
| Biology and health sciences | Perissodactyla | null |
802149 | https://en.wikipedia.org/wiki/Indian%20rhinoceros | Indian rhinoceros | The Indian rhinoceros (Rhinoceros unicornis) is also known as the greater one-horned rhinoceros, great Indian rhinoceros and Indian rhino. It is the second largest living rhinoceros species, with adult males weighing and adult females . Its thick skin is grey-brown with pinkish skin folds. It has a single horn on its snout that grows up to long. Its upper legs and shoulders are covered in wart-like bumps, and it is nearly hairless aside from the eyelashes, ear fringes and tail brush.
The Indian rhinoceros is native to the Indo-Gangetic Plain and occurs in 12 protected areas in northern India and southern Nepal. It is a grazer, eating mainly grass, but also twigs, leaves, branches, shrubs, flowers, fruits and aquatic plants. It is a largely solitary animal, only associating in the breeding season and when rearing calves. Females give birth to a single calf after a gestation of 15.7 months. The birth interval is 34–51 months. Captive individuals can live up to 47 years. It is susceptible to diseases such as anthrax, and those caused by parasites such as leeches, ticks and nematodes.
The Indian rhinoceros is listed as Vulnerable on the IUCN Red List, as the population is fragmented and restricted to less than . Excessive hunting and agricultural development reduced its range drastically. In the early 1990s, the global population was estimated at between 1,870 and 1,895 individuals. Since then, the population increased due to conservation measures taken by the governments. As of August 2018, it was estimated to comprise 3,588 individuals. However, poaching remains a continuous threat.
Taxonomy
Rhinoceros unicornis was the scientific name used by Carl Linnaeus in 1758 who described a rhinoceros with one horn. As type locality, he indicated Africa and India. He described two species in India, the other being Rhinoceros bicornis, and stated that the Indian species had two horns, while the African species had only one.
The Indian rhinoceros is a single species. Several specimens were described since the end of the 18th century under different scientific names, which are all considered synonyms of Rhinoceros unicornis today:
R. indicus by Cuvier, 1817
R. asiaticus by Blumenbach, 1830
R. stenocephalus by Gray, 1867
R. jamrachi by Sclater, 1876
R. bengalensis by Kourist, 1970
Etymology
The generic name rhinoceros is derived through Latin from the , which is composed of (rhino-, "of the nose") and (keras, "horn") with a horn on the nose. The name has been in use since the 14th century. The Latin word ūnicornis means "one-horned".
Evolution
Ancestral rhinoceroses first diverged from other perissodactyls in the Early Eocene. Mitochondrial DNA comparison suggests the ancestors of modern rhinos split from the ancestors of Equidae around 50 million years ago. The extant family, the Rhinocerotidae, first appeared in the Late Eocene in Eurasia, and the ancestors of the extant rhino species dispersed from Asia beginning in the Miocene. The last common ancestor of living rhinoceroses belonging to the subfamily Rhinocerotinae is suggested to have lived around 16 million years ago, with the ancestors of the genus Rhinoceros diverging from the ancestors of other living rhinoceroses around 15 million years ago. The genus Rhinoceros has been found to be overall slightly more closely related to the Sumatran rhinoceros (as well as to the extinct woolly rhinoceros and the extinct Eurasian genus Stephanorhinus) than to living African rhinoceroses, though there appears to have been gene flow between the ancestors of living African rhinoceroses and the genus Rhinoceros, as well as between the ancestors of the genus Rhinoceros and the ancestors of the woolly rhinoceros and Stephanorhinus.
A cladogram showing the relationships of recent and Late Pleistocene rhinoceros species (minus Stephanorhinus hemitoechus) based on whole nuclear genomes, after Liu et al., 2021:
The earliest fossils of the genus Rhinoceros date to the Late Miocene, around 8–9 million years ago. The divergence between the Indian and Javan rhinoceros is estimated to have occurred around 4.3 million years ago. The earliest representatives of the modern Indian rhinoceros appeared during the Early Pleistocene (2.6-0.8 million years ago). Fossils indicate that the Indian rhinoceros during the Pleistocene also inhabited areas considerably further east of its current distribution, including mainland Southeast Asia, South China and the island of Java, Indonesia.
Characteristics
Indian rhinos have a thick grey-brown skin with pinkish skin folds and one horn on their snout. Their upper legs and shoulders are covered in wart-like bumps. They have very little body hair, aside from eyelashes, ear fringes and tail brush. Bulls have huge neck folds. The skull is heavy with a basal length above and an occiput above . The nasal horn is slightly back-curved with a base of about by that rapidly narrows until a smooth, even stem part begins about above base. In captive animals, the horn is frequently worn down to a thick knob. The Indian rhino's single horn is present in both bulls and cows, but not on newborn calves. The horn is pure keratin, like human fingernails, and starts to show after about six years. In most adults, the horn reaches a length of about , but has been recorded up to in length and in weight.
Among terrestrial land mammals native to Asia, Indian rhinos are second in size only to the Asian elephant. They are also the second-largest living rhinoceros, behind only the white rhinoceros. Bulls have a head and body length of with a shoulder height of , while cows have a head and body length of and a shoulder height of . The bull, averaging about is heavier than the cow, at an average of about . The largest individuals reportedly weigh up to .
The rich presence of blood vessels underneath the tissues in folds gives them the pinkish colour. The folds in the skin increase the surface area and help in regulating the body temperature. The thick skin does not protect against bloodsucking Tabanus flies, leeches and ticks.
Distribution and habitat
Indian rhinos once ranged across the entire northern part of the Indian subcontinent, along the Indus, Ganges and Brahmaputra River basins, from Pakistan to the Indian-Myanmar border, including Bangladesh and the southern parts of Nepal and Bhutan. They may have also occurred in Myanmar, southern China and Indochina. They inhabit the alluvial grasslands of the Terai and the Brahmaputra basin. As a result of habitat destruction and climatic changes its range has gradually been reduced so that by the 19th century, it only survived in the Terai grasslands of southern Nepal, northern Uttar Pradesh, northern Bihar, northern West Bengal, and in the Brahmaputra Valley of Assam.
The species was present in northern Bihar and Oudh at least until 1770 as indicated in maps produced by Colonel Gentil. On the former abundance of the species, Thomas C. Jerdon wrote in 1867:
Today, its range has further shrunk to a few pockets in southern Nepal, northern West Bengal, and the Brahmaputra Valley. Its habitat is surrounded by human-dominated landscapes, so that in many areas, it occurs in cultivated areas, pastures, and secondary forests. In the 1980s, Indian rhinos were frequently seen in the narrow plain area of Manas River and Royal Manas National Park in Bhutan.
Populations
In 2022, the total Indian rhinoceros population was estimated to be 4,014 individuals, up from 2,577 in 2006. Among them, 3,262 are in India and the remaining 752 are in Nepal and Bhutan. There is no permanent rhino population in Bhutan, but small rhino populations are occasionally known to cross from the Manas National Park or Buxa Tiger Reserve in India.
In India, there are around 2,885 individuals in Assam, including 2,613 in Kaziranga National Park, 125 in Orang National Park, 107 in Pobitora Wildlife Sanctuary and 40 in Manas National Park. West Bengal has a population of 339 individuals, including 287 in Jaldapara National Park and 52 in Gorumara National Park. Only 38 individuals are found in Dudhwa National Park, in Uttar Pradesh.
By 2014, the population in Assam increased to 2,544 Indian rhinos, an increase of 27% since 2006, although more than 150 individuals were killed by poachers during these years.
The population in Kaziranga National Park was estimated at 2,048 individuals in 2009. By 2009, the population in Pobitora Wildlife Sanctuary had increased to 84 individuals in an area of .
In 2015, Nepal had 645 Indian rhinos living in Parsa National Park, Chitwan National Park, Bardia National Park, Shuklaphanta Wildlife Reserve and respective buffer zones in the Terai Arc Landscape as recorded in a survey conducted from 11 April to 2 May 2015. The survey showed that the population of rhinos in Nepal from 2011 to 2015 increased 21% or 111 individuals.
The Indian rhino population, which once numbered as low as 100 individuals in the early 1900s, has increased to more than 3,700 in the year 2021 as per The International Rhino Foundation.
Ecology and behaviour
Bulls are usually solitary. Groups consist of cows with calves, or of up to six subadults. Such groups congregate at wallows and grazing areas. They are foremost active in early mornings, late afternoons and at night, but rest during hot days.
They bathe regularly. The folds in their skin trap water and hold it even when they exit wallows. They are excellent swimmers and can run at speeds of up to for short periods. They have excellent senses of hearing and smell, but relatively poor eyesight. Over 10 distinct vocalisations have been recorded. Males have home ranges of around that overlap each other. Dominant males tolerate other males passing through their territories except when they are in mating season, when dangerous fights break out.
Indian rhinos have few natural enemies, except for tigers, which sometimes kill unguarded calves, but adult rhinos are less vulnerable due to their size. Mynahs and egrets both eat invertebrates from the rhino's skin and around its feet. Tabanus flies, a type of horse-fly, are known to bite rhinos. The rhinos are also vulnerable to diseases spread by parasites such as leeches, ticks, and nematodes like Bivitellobilharzia nairi. Anthrax and the blood-disease sepsis are known to occur.
In March 2017, a group of four tigers consisting of an adult male, tigress and two cubs killed a 20-year-old male Indian rhinoceros in Dudhwa Tiger Reserve. Such cases are rare, as Indian rhinoceroseslike most megaherbivoresare mostly invulnerable to predation.
Diet
Indian rhinos are grazers. Their diet consists almost entirely of grasses (such as Arundo donax, Bambusa tulda, Cynodon dactylon, and Oryza sativa), but they also eat leaves, twigs and branches of shrubs and trees (such as Lagerstroemia indica), flowers, fruits (such as Ficus religiosa), and submerged and floating aquatic plants.
They feed in the mornings and evenings. They use their semi-prehensile lips to grasp grass stems, bend the stem down, bite off the top, and then eat the grass. They tackle very tall grasses or saplings by walking over the plant, with legs on both sides and using the weight of their bodies to push the end of the plant down to the level of the mouth. Mothers also use this technique to make food edible for their calves. They drink for a minute or two at a time, often imbibing water filled with rhinoceros urine.
Social life
Indian rhinos form a variety of social groupings. Bulls are generally solitary, except for mating and fighting. Cows are largely solitary when they are without calves. Mothers will stay close to their calves for up to four years after their birth, sometimes allowing an older calf to continue to accompany her once a newborn calf arrives. Subadult bulls and cows form consistent groupings, as well. Groups of two or three young bulls often form on the edge of the home ranges of dominant bulls, presumably for protection in numbers. Young cows are slightly less social than the bulls. Indian rhinos also form short-term groupings, particularly at forest wallows during the monsoon season and in grasslands during March and April. Groups of up to 10 rhinos, typically a dominant male with females and calves, gather in wallows.
Indian rhinos make a wide variety of vocalisations. At least 10 distinct vocalisations have been identified: snorting, honking, bleating, roaring, squeak-panting, moo-grunting, shrieking, groaning, rumbling and humphing. In addition to noises, the Indian rhino uses olfactory communication. Adult bulls urinate backwards, as far as behind them, often in response to being disturbed by observers. Like all rhinos, the Indian rhinoceros often defecates near other large dung piles. The Indian rhino has pedal scent glands which are used to mark their presence at these rhino latrines. Bulls have been observed walking with their heads to the ground as if sniffing, presumably following the scent of cows.
In aggregations, Indian rhinos are often friendly. They will often greet each other by waving or bobbing their heads, mounting flanks, nuzzling noses, or licking. Indian rhinos will playfully spar, run around, and play with twigs in their mouths. Adult bulls are the primary instigators in fights. Fights between dominant bulls are the most common cause of rhino mortality, and bulls are also very aggressive toward cows during courtship. Bulls chase cows over long distances and even attack them face-to-face. Indian rhinos use their horns for fighting, albeit less frequently than African rhinos, largely using the incisors of the lower jaw to inflict wounds.
Reproduction
Captive bulls breed at five years of age, but wild bulls attain dominance much later when they are larger. In one five-year field study, only one Indian rhino estimated to be younger than 15 years mated successfully. Captive cows breed as young as four years of age, but in the wild, they usually start breeding only when six years old, which likely indicates they need to be large enough to avoid being killed by aggressive bulls. The ovarian cycle lasts 5.5 to 9 weeks on average. Their gestation period is around 15.7 months, and birth interval ranges from 34 to 51 months. An estimated 10% of calves will die before maturity. This is mainly attributed to predatory attacks from tigers (Panthera tigris). In captivity, four Indian rhinos lived over 40 years, the oldest living to be 47.
Threats
Habitat degradation caused by human activities and climate change as well as the resulting increase in the floods has caused many Indian rhino deaths and has limited their ranging areas which is shrinking.
Serious declines in quality of habitat have occurred in some areas, due to severe invasion by alien plants into grasslands affecting some populations, and demonstrated reductions in the extent of grasslands and wetland habitats due to woodland encroachment and silting up of beels (swampy wetlands). Grazing by domestic livestock is another cause.
The Indian rhino species is inherently at risk because over 70% of its population occurs at a single site, Kaziranga National Park. Any catastrophic event such as disease, civil disorder, poaching, or habitat loss would have a devastating impact on the Indian rhino's status. Additionally, a small population of rhinos may be prone to inbreeding depression.
Poaching
Sport hunting became common in the late 19th and early 20th centuries and was the main cause for the decline of Indian rhinoceros populations. Indian rhinos were hunted relentlessly and persistently. Reports from the mid-19th century claim that some British military officers shot more than 200 rhinos in Assam alone. By 1908, the population in Kaziranga National Park had decreased to around 12 individuals. In the early 1900s, the Indian rhinoceros was almost extinct. At present, poaching for the use of horn in traditional Chinese Medicine is one of the main threats that has led to decreases in several important populations. Poaching for the Indian rhino's horn became the single most important reason for the decline of the Indian rhinoceros after conservation measures were put in place from the beginning of the 20th century, when legal hunting ended. From 1980 to 1993, 692 rhinos were poached in India, including 41 rhinos in India's Laokhowa Wildlife Sanctuary in 1983, almost the entire population of the sanctuary. By the mid-1990s, the Indian rhinoceros had been extirpated in this sanctuary. Between 2000 and 2006, more than 150 rhinos were poached in Assam. Almost 100 rhinos were poached in India between 2013 and 2018.
In 1950, in Nepal the Chitwan's forest and grasslands extended over more than and were home to about 800 rhinos. When poor farmers from the mid-hills moved to the Chitwan Valley in search of arable land, the area was subsequently opened for settlement, and poaching of wildlife became rampant. The Chitwan population has repeatedly been jeopardised by poaching; in 2002 alone, poachers killed 37 animals to saw off and sell their valuable horns.
Conservation
The Indian rhinoceros is listed as vulnerable by the IUCN Red list, as of 2018. Globally, R. unicornis has been listed in CITES Appendix I since 1975. The Indian and Nepalese governments have taken major steps towards Indian rhinoceros conservation, especially with the help of the World Wide Fund for Nature (WWF) and other non-governmental organisations. In 1910, all rhino hunting in India became prohibited.
In 1957, the country's first conservation law ensured the protection of rhinos and their habitat. In 1959, Edward Pritchard Gee undertook a survey of the Chitwan Valley, and recommended the creation of a protected area north of the Rapti River and of a wildlife sanctuary south of the river for a trial period of 10 years. After his subsequent survey of Chitwan in 1963, he recommended extension of the sanctuary to the south. By the end of the 1960s, only 95 rhinos remained in the Chitwan Valley. The dramatic decline of the rhino population and the extent of poaching prompted the government to institute the Gaida Gasti – a rhino reconnaissance patrol of 130 armed men and a network of guard posts all over Chitwan. To prevent the extinction of rhinos, the Chitwan National Park was gazetted in December 1970, with borders delineated the following year and established in 1973, initially encompassing an area of . To ensure the survival of rhinos in case of epidemics, animals were translocated annually from Chitwan to Bardia National Park and Shuklaphanta National Park since 1986.
The Indian rhinoceros population living in Chitwan and Parsa National Parks was estimated at 608 mature individuals in 2015.
Reintroduction to new areas
Indian rhinos have been reintroduced to areas where they had previously inhabited but became extinct. These efforts have produced mixed results, mainly due to lack of proper planning and management, sustained effort, and adequate security for the introduced animals.
In 1984, five Indian rhinos were relocated to Dudhwa National Park—four from the fields outside the Pobitora Wildlife Sanctuary and one from Goalpara. This has born results and the population has increased to 21 rhinos by 2006. In early 1980s, Laokhowa Wildlife Sanctuary in Assam had more than 70 Indian rhinos which were all killed by poachers. In 2016, two Indian rhinos, a mother and her daughter, were reintroduced to the sanctuary from Kaziranga National Park as part of the Indian Rhino Vision 2020 (IRV 2020) program, but both animals died within months due to natural causes.
Indian rhinos were once found as far west as the Peshawar Valley during the reign of Mughal Emperor Babur, but are now extinct in Pakistan. After rhinos became "regionally extinct" in Pakistan, two rhinos from Nepal were introduced in 1983 to Lal Suhanra National Park, which have not bred so far.
In captivity
Indian rhinoceroses were initially difficult to breed in captivity. In the second half of the 20th century, zoos became adept at breeding Indian rhinoceros. By 1983, nearly 40 babies had been born in captivity. As of 2012, 33 Indian rhinos were born at Switzerland's Zoo Basel alone, meaning that most captive animals are related to the Basel population. Due to the success of Zoo Basel's breeding program, the International Studbook for the species has been kept there since 1972. Since 1990, the Indian rhino European Endangered Species Programme is also being coordinated there, with the goal of maintaining genetic diversity in the global captive Indian rhinoceros population.
The first recorded captive birth of an Indian rhinoceros was in Kathmandu in 1826, but another successful birth did not occur for nearly 100 years. In 1925, a rhino was born in Kolkata. No rhinoceros was successfully bred in Europe until 1956 when first European breeding took place when baby rhino Rudra was born in Zoo Basel on 14 September 1956. In June 2009, an Indian rhino was artificially inseminated using sperm collected four years previously and cryopreserved at the Cincinnati Zoo's CryoBioBank before being thawed and used. She gave birth to a male calf in October 2010. In June 2014, the first "successful" live-birth from an artificially inseminated rhino took place at the Buffalo Zoo in New York. As in Cincinnati, cryopreserved sperm was used to produce the female calf, Monica.
Cultural significance
The Indian rhinoceros is one of the motifs on the Pashupati seal and many terracotta figurines that were excavated at archaeological sites of the Indus Valley civilisation.
The Rhinoceros Sutra is an early text in the Buddhist tradition and is part of the Gandhāran Buddhist texts and the Pali Canon; a version was also incorporated into the Sanskrit Mahavastu. It praises the solitary lifestyle and stoicism of the Indian rhinoceros and is associated with the eremitic lifestyle symbolized by the Pratyekabuddha.
Europe
In the 3rd century, Philip the Arab exhibited an Indian rhinoceros in Rome. In 1515, Manuel I of Portugal obtained an Indian rhinoceros as a gift, which he passed on to Pope Leo X, but which died in a shipwreck off the coast of Italy in early 1516, on the way from Lisbon to Rome. Three artistic representations were prepared of this rhinoceros: A woodcut by Hans Burgkmair, a drawing and a woodcut called Dürer's Rhinoceros by Albrecht Dürer, all dated 1515. In 1577–1588, Abada was a female Indian rhinoceros kept by the Portuguese kings Sebastian I and Henry I from 1577 to 1580 and by Philip II of Spain from about 1580 to 1588. She was the first rhinoceros seen in Europe after Dürer's Rhinoceros. In about 1684, the first presumably Indian rhinoceros arrived in England. George Jeffreys, 1st Baron Jeffreys spread the rumour that his chief rival Francis North, 1st Baron Guilford had been seen riding on it.
In 1741–1758, Clara the rhinoceros (c. 1738 – 14 April 1758) was a female Indian rhinoceros who became famous during 17 years of touring Europe in the mid-18th century. She arrived in Europe in Rotterdam in 1741, becoming the fifth living rhinoceros to be seen in Europe in modern times since Dürer's rhinoceros in 1515. After tours through towns in the Dutch Republic, the Holy Roman Empire, Switzerland, the Polish–Lithuanian Commonwealth, France, the Kingdom of the Two Sicilies, the Papal States, Bohemia and Denmark, she died in Lambeth, England. In 1739, she was drawn and engraved by two English artists. She was then brought to Amsterdam, where Jan Wandelaar made two engravings that were published in 1747. In the subsequent years, the rhinoceros was exhibited in several European cities. In 1748, Johann Elias Ridinger made an etching of her in Augsburg, and Petrus Camper modelled her in clay in Leiden. In 1749, Georges-Louis Leclerc, Comte de Buffon drew it in Paris. In 1751, Pietro Longhi painted her in Venice.
| Biology and health sciences | Perissodactyla | Animals |
12895973 | https://en.wikipedia.org/wiki/Common%20walkingstick | Common walkingstick | The common walkingstick or northern walkingstick (Diapheromera femorata) is a species of phasmid or stick insect found across North America. The average length of this species is 75mm (3 in) for males and 95mm (3.7 in) for females.
The insect is found in deciduous forest throughout North America, where it eats many types of plant foliage. Even though the common walkingstick is a generalist it does tend to prefer foliage from oak and hazelnut trees. Localised clusters of these insects sometimes occur; however, the insects have no wings, and dispersal from tree to tree is limited.
Description
The common walkingstick is a slender, elongated insect that camouflages itself by resembling a twig. The sexes differ, with the male usually being brown and about in length while the female is greenish-brown, and rather larger at . There are three pairs of legs, but at rest, the front pair is extended forward beside the antennae, forming an extension of the twig-like effect. Neither sex has wings, the antennae are two thirds of the length of the body, and each of the cerci (paired appendages at the tip of the abdomen) has a single segment.
Distribution
The walkingstick is native to North America. Its range extends from the Atlantic coast from Maine to Florida, as far west as California and northwards to North Dakota. It also occurs in Canada (where it is the only stick insect) being present in Alberta, Manitoba, Ontario and Québec.
Ecology
D. femorata is herbivorous, feeding mainly on the leaves of trees. They are leaf skeletonisers, eating the tissues between the leaf veins, pausing for a while and then walking on to new leaves. They can feed at any time of day but the greatest feeding activity has been noted between 9pm and 3am. Early-stage nymphs are often found on American hazel (Corylus americana) and black cherry (Prunus serotina), but where these are scarce, they are likely to be on white oak (Quercus alba). Older individuals may choose black oak (Quercus velutina). Another food tree is the black locust (Robinia pseudoacacia). Adults are present in August and September in the northern part of the range, but because of their tendency to feed high in the canopy, the insects are seldom seen.
The stick insect life cycle is hemimetabolous, proceeding through a series of nymphal stages. Breeding takes place in late summer after the nymph has moulted for the last time and become an adult. Egg-laying takes place about a week after mating and the eggs, measuring across, are dropped singly to the forest floor. Here they overwinter in the leaf litter, hatching the following May or even a year later. If conditions are dry, the newly hatched young may fail to extricate themselves from their egg-capsules. Ones that succeed in doing so climb up the trunks of trees, start to feed on the foliage, and pass through four to six moults as they grow.
| Biology and health sciences | Insects: General | Animals |
1258615 | https://en.wikipedia.org/wiki/Volcanic%20crater | Volcanic crater | A volcanic crater is an approximately circular depression in the ground caused by volcanic activity. It is typically a bowl-shaped feature containing one or more vents. During volcanic eruptions, molten magma and volcanic gases rise from an underground magma chamber, through a conduit, until they reach the crater's vent, from where the gases escape into the atmosphere and the magma is erupted as lava. A volcanic crater can be of large dimensions, and sometimes of great depth. During certain types of explosive eruptions, a volcano's magma chamber may empty enough for an area above it to subside, forming a type of larger depression known as a caldera.
Geomorphology
In most volcanoes, the crater is situated at the top of a mountain formed from the erupted volcanic deposits such as lava flows and tephra. Volcanoes that terminate in such a summit crater are usually of a conical form. Other volcanic craters may be found on the flanks of volcanoes, and these are commonly referred to as flank craters. Some volcanic craters may fill either fully or partially with rain and/or melted snow, forming a crater lake.
These lakes may become soda lakes, many of which are associated with active tectonic and volcanic zones.
A crater may be breached during an eruption, either by explosions or by lava, or through later erosion. Breached craters have a much lower rim on one side.
Some volcanoes, such as maars, consist of a crater alone, with scarcely any mountain at all. These volcanic explosion craters are formed when magma rises through water-saturated rocks, which causes a phreatic eruption. Volcanic craters from phreatic eruptions often occur on plains away from other obvious volcanoes. Not all volcanoes form craters.
| Physical sciences | Volcanology | Earth science |
1258972 | https://en.wikipedia.org/wiki/Evening | Evening | Evening is the period of a day that begins at the end of daylight and overlaps with the beginning of night. It generally indicates the period of time when the sun is close to the horizon and comprises the periods of civil, nautical and astronomical twilight. The exact times when evening begins and ends are subjective and depend on location and time of year. It may be used colloquially to include the last waning daytime shortly before sunset.
Etymology
The word is derived from the Old English ǣfnung, meaning 'the coming of evening, sunset, time around sunset', which originated from æfnian, meaning "become evening, grow toward evening". The Old English æfnian originated from æfen (eve), which meant "the time between sunset and darkness", and was synonymous with even (Old English æfen), which meant the end of the day. The use of "evening" dates from the mid 15th century.
Start time
The Encyclopædia Britannica defines evening as varying according to daylight and lifestyle, but says that many people consider it to begin at 5 p.m.
In a social context, the Oxford English Dictionary defines evening as "the time from about 6 p.m., or sunset if earlier".
As such, there is no fixed consensus on when the period of evening starts.
Astronomy
Despite the less favorable lighting conditions for optical astronomy, evening can be useful for observing objects orbiting close to the Sun. Evening (and morning) serves as the optimum time for viewing the inferior planets Venus and Mercury. It is a popular time to hunt for comets, as their tails grow more prominent as these objects draw closer to the Sun. The evening (and morning) twilight is used to search for near-Earth asteroids that orbit inside the orbit of the Earth. In mid-latitudes, spring evenings around the time of the equinox―that is, the March one in the Northern Hemisphere and the September equinox to the south of the equator―are favorable for viewing the zodiacal light.
| Physical sciences | Celestial mechanics | Astronomy |
1259241 | https://en.wikipedia.org/wiki/Habitat%20fragmentation | Habitat fragmentation | Habitat fragmentation describes the emergence of discontinuities (fragmentation) in an organism's preferred environment (habitat), causing population fragmentation and ecosystem decay. Causes of habitat fragmentation include geological processes that slowly alter the layout of the physical environment (suspected of being one of the major causes of speciation), and human activity such as land conversion, which can alter the environment much faster and causes the extinction of many species. More specifically, habitat fragmentation is a process by which large and contiguous habitats get divided into smaller, isolated patches of habitats.
Definition
The term habitat fragmentation includes five discrete phenomena:
Reduction in the total area of the habitat
Decrease of the interior: edge ratio
Isolation of one habitat fragment from other areas of habitat
Breaking up of one patch of habitat into several smaller patches
Decrease in the average size of each patch of habitat
"fragmentation ... not only causes loss of the amount of habitat but by creating small, isolated patches it also changes the properties of the remaining habitat" (van den Berg et al. 2001). Habitat fragmentation is the landscape level of the phenomenon, and patch level process. Thus meaning, it covers; the patch areas, edge effects, and patch shape complexity.
In scientific literature, there is some debate whether the term "habitat fragmentation" applies in cases of habitat loss, or whether the term primarily applies to the phenomenon of habitat being cut into smaller pieces without significant reduction in habitat area. Scientists who use the stricter definition of "habitat fragmentation" per se would refer to the loss of habitat area as "habitat loss" and explicitly mention both terms if describing a situation where the habitat becomes less connected and there is less overall habitat.
Furthermore, habitat fragmentation is considered as an invasive threat to biodiversity, due to its implications of affecting large number of species than biological invasions, overexploitation, or pollution.
Additionally, the effects of habitat fragmentation damage the ability for species, such as native plants, to be able to effectively adapt to their changing environments. Ultimately, this prevents gene flow from one generation of population to the next, especially for species living in smaller population sizes. Whereas, for species of larger populations have more genetic mutations which can arise and genetic recombination impacts which can increase species survival in those environments. Overall, habitat fragmentation results in habitat disintegration and habitat loss which both tie into destructing biodiversity as a whole.
Causes
Natural causes
Evidence of habitat destruction through natural processes such as volcanism, fire, and climate change is found in the fossil record.Studies have demonstrated the impacts of individual species at the landscape level For example, From research the results show that the impact of deer herbivory on forest plant communities can be observed at the landscape level at the Rondeau Provincial park for the period of 1955-1978and also, habitat fragmentation of tropical rainforests in Euramerica 300 million years ago led to a great loss of amphibian diversity, but simultaneously the drier climate spurred on a burst of diversity among reptiles.
Human causes
Habitat fragmentation is frequently caused by humans when native plants are cleared for human activities such as agriculture, rural development, urbanization and the creation of hydroelectric reservoirs. Habitats which were once continuous become divided into separate fragments. Due to human activities, many tropical and temperate habitats have already been severely fragmented, and in the near future, the degree of fragmentation will significantly rise. After intensive clearing, the separate fragments tend to be very small islands isolated from each other by cropland, pasture, pavement, or even barren land. The latter is often the result of slash and burn farming in tropical forests. In the wheat belt of central-western New South Wales, Australia, 90% of the native vegetation has been cleared and over 99% of the tall grass prairie of North America has been cleared, resulting in extreme habitat fragmentation.
Endogenous vs. exogenous
There are two types of processes that can lead to habitat fragmentation. There are exogenous processes and endogenous processes. Endogenous is a process that develops as a part of species biology so they typically include changes in biology, behavior, and interactions within or between species. Endogenous threats can result in changes to breeding patterns or migration patterns and are often triggered by exogenous processes. Exogenous processes are independent of species biology and can include habitat degradation, habitat subdivision or habitat isolation. These processes can have a substantial impact on endogenous processes by fundamentally altering species behavior. Habitat subdivision or isolation can lead to changes in dispersal or movement of species including changes to seasonal migration. These changes can lead to a decrease in a density of species, increased competition or even increased predation.
Implications
Habitat and biodiversity loss
One of the major ways that habitat fragmentation affects biodiversity is by reducing the amount of suitable habitat available for organisms. Habitat fragmentation often involves both habitat destruction and the subdivision of previously continuous habitat. Plants and other sessile organisms are disproportionately affected by some types of habitat fragmentation because they cannot respond quickly to the altered spatial configuration of the habitat. Habitat fragmentation consistently reduces biodiversity by 13 to 75% and impairs key ecosystem functions by decreasing biomass and altering nutrient cycles. This underscores the severe and lasting ecological impacts of fragmentation, which could be highlighted in the sections discussing the consequences of fragmentation.
Habitat loss, which can occur through the process of habitat fragmentation, is considered to be the greatest threat to species. But, the effect of the configuration of habitat patches within the landscape, independent of the effect of the amount of habitat within the landscape (referred to as fragmentation per se), has been suggested to be small. A review of empirical studies found that, of the 381 reported significant effect of habitat fragmentation per se on species occurrences, abundances or diversity in the scientific literature, 76% were positive whereas 24% were negative. Despite these results, the scientific literature tends to emphasize negative effects more than positive effects. Positive effects of habitat fragmentation per se imply that several small patches of habitat can have higher conservation value than a single large patch of equivalent size. Land sharing strategies could therefore have more positive impacts on species than land sparing strategies. Although the negative effects of habitat loss are generally viewed to be much larger than that of habitat fragmentation, the two events are heavily connected and observations are not usually independent of one another.
Area is the primary determinant of the number of species in a fragment and the relative contributions of demographic and genetic processes to the risk of global population extinction depend on habitat configuration, stochastic environmental variation and species features. Minor fluctuations in climate, resources, or other factors that would be unremarkable and quickly corrected in large populations can be catastrophic in small, isolated populations. Thus fragmentation of habitat is an important cause of species extinction. Population dynamics of subdivided populations tend to vary asynchronously. In an unfragmented landscape a declining population can be "rescued" by immigration from a nearby expanding population. In fragmented landscapes, the distance between fragments may prevent this from happening. Additionally, unoccupied fragments of habitat that are separated from a source of immigrants by some barrier are less likely to be repopulated than adjoining fragments. Even small species such as the Columbia spotted frog are reliant on the rescue effect. Studies showed 25% of juveniles travel a distance over 200m compared to 4% of adults. Of these, 95% remain in their new locale, demonstrating that this journey is necessary for survival.
Additionally, habitat fragmentation leads to edge effects. Microclimatic changes in light, temperature, and wind can alter the ecology around the fragment, and in the interior and exterior portions of the fragment. Fires become more likely in the area as humidity drops and temperature and wind levels rise. Exotic and pest species may establish themselves easily in such disturbed environments, and the proximity of domestic animals often upsets the natural ecology. Also, habitat along the edge of a fragment has a different climate and favours different species from the interior habitat. Small fragments are therefore unfavourable for species that require interior habitat. The percentage preservation of contiguous habitats is closely related to both genetic and species biodiversity preservation. Generally a 10% remnant contiguous habitat will result in a 50% biodiversity loss.
Much of the remaining terrestrial wildlife habitat in many third world countries has experienced fragmentation through the development of urban expansion such as roads interfering with habitat loss. Aquatic species’ habitats have been fragmented by dams and water diversions. These fragments of habitat may not be large or connected enough to support species that need a large territory where they can find mates and food. The loss and fragmentation of habitats makes it difficult for migratory species to find places to rest and feed along their migration routes.
The effects of current fragmentation will continue to emerge for decades. Extinction debts are likely to come due, although the counteracting immigration debts may never fully be paid. Indeed, the experiments here reveal ongoing losses of biodiversity and ecosystem functioning two decades or longer after fragmentation occurred. Understanding the relationship between transient and long-term dynamics is a substantial challenge that ecologists must tackle, and fragmentation experiments will be central for relating observation to theory.
Informed conservation
Habitat fragmentation is often a cause of species becoming threatened or endangered. The existence of viable habitat is critical to the survival of any species, and in many cases, the fragmentation of any remaining habitat can lead to difficult decisions for conservation biologists. Given a limited amount of resources available for conservation is it preferable to protect the existing isolated patches of habitat or to buy back land to get the largest possible contiguous piece of land. In rare cases, a conservation reliant species may gain some measure of disease protection by being distributed in isolated habitats, and when controlled for overall habitat loss some studies have shown a positive relationship between species richness and fragmentation; this phenomenon has been called the habitat amount hypothesis, though the validity of this claim has been disputed. The ongoing debate of what size fragments are most relevant for conservation is often referred to as SLOSS (Single Large or Several Small). Habitat loss in a biodiversity hotspot can result in a localized extinction crisis, generally speaking habitat loss in a hotspot location can be a good indicator or predictor of the number of threatened and extinct endemic species.
One solution to the problem of habitat fragmentation is to link the fragments by preserving or planting corridors of native vegetation. In some cases, a bridge or underpass may be enough to join two fragments. This has the potential to mitigate the problem of isolation but not the loss of interior habitat. Wildlife corridors can help animals to move and occupy new areas when food sources or other natural resources are lacking in their core habitat, and animals can find new mates in neighbouring regions so that genetic diversity can increase. Species that relocate seasonally can do so more safely and effectively when it does not interfere with human development barriers.
Due to the continuous expansion of urban landscapes, current research is looking at green roofs being possible vectors of habitat corridors. A recent study has found that green roofs are beneficial in connecting the habitats of arthropods, specifically bees and weevils.
Another mitigation measure is the enlargement of small remnants to increase the amount of interior habitat. This may be impractical since developed land is often more expensive and could require significant time and effort to restore.
The best solution is generally dependent on the particular species or ecosystem that is being considered. More mobile species, like most birds, do not need connected habitat while some smaller animals, like rodents, may be more exposed to predation in open land. These questions generally fall under the headings of metapopulations island biogeography.
Genetic risks
As the remaining habitat patches are smaller, they tend to support smaller populations of fewer species. Small populations are at an increased risk of a variety of genetic consequences that influence their long-term survival. Remnant populations often contain only a subset of the genetic diversity found in the previously continuous habitat. In these cases, processes that act upon underlying genetic diversity, such as adaptation, have a smaller pool of fitness-maintaining alleles to survive in the face of environmental change. However, in some scenarios, where subsets of genetic diversity are partitioned among multiple habitat fragments, almost all original genetic diversity can be maintained despite each individual fragment displaying a reduced subset of diversity.
Gene Flow and Inbreeding
Gene flow occurs when individuals of the same species exchange genetic information through reproduction. Populations can maintain genetic diversity through migration. When a habitat becomes fragmented and reduced in area, gene flow and migration are typically reduced. Fewer individuals will migrate into the remaining fragments, and small disconnected populations that may have once been part of a single large population will become reproductively isolated. Scientific evidence that gene flow is reduced due to fragmentation depends on the study species. While trees that have long-range pollination and dispersal mechanisms may not experience reduced gene flow following fragmentation, most species are at risk of reduced gene flow following habitat fragmentation.
Reduced gene flow, and reproductive isolation can result in inbreeding between related individuals. Inbreeding does not always result in negative fitness consequences, but when inbreeding is associated with fitness reduction it is called inbreeding depression. Inbreeding becomes of increasing concern as the level of homozygosity increases, facilitating the expression of deleterious alleles that reduce the fitness. Habitat fragmentation can lead to inbreeding depression for many species due to reduced gene flow. Inbreeding depression is associated with conservation risks, like local extinction.
Genetic drift
Small populations are more susceptible to genetic drift. Genetic drift is random changes to the genetic makeup of populations and leads to reductions in genetic diversity. The smaller the population is, the more likely genetic drift will be a driving force of evolution rather than natural selection. Because genetic drift is a random process, it does not allow species to become more adapted to their environment. Habitat fragmentation is associated with increases to genetic drift in small populations which can have negative consequences for the genetic diversity of the populations. However, research suggests that some tree species may be resilient to the negative consequences of genetic drift until population size is as small as ten individuals or less.
Genetic consequences of habitat fragmentation for plant populations
Habitat fragmentation decreases the size and increases plant populations' spatial isolation. With genetic variation and increased methods of inter-population genetic divergence due to increased effects of random genetic drift, elevating inbreeding and reducing gene flow within plant species. While genetic variation may decrease with remnant population size, not all fragmentation events lead to genetic losses and different types of genetic variation. Rarely, fragmentation can also increase gene flow among remnant populations, breaking down local genetic structure.
Adaptation
In order for populations to evolve in response to natural selection, they must be large enough that natural selection is a stronger evolutionary force than genetic drift. Recent studies on the impacts of habitat fragmentation on adaptation in some plant species have suggested that organisms in fragmented landscapes may be able to adapt to fragmentation. However, there are also many cases where fragmentation reduces adaptation capacity because of small population size.
Examples of impacted species
Some species that have experienced genetic consequences due to habitat fragmentation are listed below:
Macquaria australasica
Fagus sylvatica
Betula nana
Rhinella ornata
Ochotona princeps
Uta stansburiana
Plestiodon skiltonianus
Sceloporus occidentalis
Chamaea fasciata
Effect on animal behaviours
Although the way habitat fragmentation affects the genetics and extinction rates of species has been heavily studied, fragmentation has also been shown to affect species' behaviours and cultures as well. This is important because social interactions can determine and have an effect on a species' fitness and survival. Habitat fragmentation alters the resources available and the structure of habitats, as a result, alters the behaviours of species and the dynamics between differing species. Behaviours affected can be within a species such as reproduction, mating, foraging, species dispersal, communication and movement patterns or can be behaviours between species such as predator-prey relationships. In addition, when animals happen to venture into unknown areas in between fragmented forests or landscapes, they can supposedly come into contact with humans which puts them at a great risk and further decreases their chances of survival.
Predation behaviours
Habitat fragmentation due to anthropogenic activities has been shown to greatly affect the predator-prey dynamics of many species by altering the number of species and the members of those species. This affects the natural predator-prey relationships between animals in a given community and forces them to alter their behaviours and interactions, therefore resetting the so-called "behavioral space race". The way in which fragmentation changes and re-shapes these interactions can occur in many different forms. Most prey species have patches of land that are a refuge from their predators, allowing them the safety to reproduce and raise their young. Human introduced structures such as roads and pipelines alter these areas by facilitating predator activity in these refuges, increasing predator-prey overlap. The opposite could also occur in the favour of prey, increasing prey refuge and subsequently decreasing predation rates. Fragmentation may also increase predator abundance or predator efficiency and therefore increase predation rates in this manner. Several other factors can also increase or decrease the extent to which the shifting predator-prey dynamics affect certain species, including how diverse a predators diet is and how flexible habitat requirements are for predators and prey. Depending on which species are affected and these other factors, fragmentation and its effects on predator-prey dynamics may contribute to species extinction. In response to these new environmental pressures, new adaptive behaviours may be developed. Prey species may adapt to increased risk of predation with strategies such as altering mating tactics or changing behaviours and activities related to food and foraging.
Boreal woodland caribous
In the boreal woodland caribous of British Columbia, the effects of fragmentation are demonstrated. The species refuge area is peatland bog which has been interrupted by linear features such as roads and pipelines. These features have allowed their natural predators, the wolf, and the black bear to more efficiently travel over landscapes and between patches of land. Since their predators can more easily access the caribous' refuge, the females of the species attempt to avoid the area, affecting their reproductive behaviours and offspring produced.
Communication behaviours
Fragmentation affecting the communication behaviours of birds has been well studied in Dupont's Lark. The Larks primarily reside in regions of Spain and are a small passerine bird which uses songs as a means of cultural transmission between members of the species. The Larks have two distinct vocalizations, the song, and the territorial call. The territorial call is used by males to defend and signal territory from other male Larks and is shared between neighbouring territories when males respond to a rivals song. Occasionally it is used as a threat signal to signify an impending attack on territory. A large song repertoire can enhance a male's ability to survive and reproduce as he has a greater ability to defend his territory from other males, and a larger number of males in the species means a larger variety of songs being transmitted. Fragmentation of the Dupont's Lark territory from agriculture, forestry and urbanization appears to have a large effect on their communication structures. Males only perceive territories of a certain distance to be rivals and so isolation of territory from others due to fragmentation leads to a decrease in territorial calls as the males no longer have any reason to use it or have any songs to match.
Humans have also brought on varying implications into ecosystems which in turn affect animal behaviour and responses generated. Although there are some species which are able to survive these kinds of harsh conditions, such as, cutting down wood in the forests for pulp and paper industries, there are animals which can survive this change but some that cannot. An example includes, varying aquatic insects are able to identify appropriate ponds to lay their eggs with the aid of polarized light to guide them, however, due to ecosystem modifications caused by humans they are led onto artificial structures which emit artificial light which are induced by dry asphalt dry roads for an example.
Effect on microorganisms
While habitat fragmentation is often associated with its effects on large plant and animal populations and biodiversity, due to the interconnectedness of ecosystems there are also significant effects that it has on the microbiota of an environment. Increased fragmentation has been linked to reduced populations and diversity of fungi responsible for decomposition, as well as the insects they are host to. This has been linked to simplified food webs in highly fragmented areas compared to old growth forests. Furthermore, edge effects have been shown to result in significantly varied microenvironments compared to interior forest due to variations in light availability, presence of wind, changes in precipitation, and overall moisture content of leaf litter. These microenvironments are often not conducive to overall forest health as they enable generalist species to thrive at the expense of specialists that depend on specific environments.
Effect on mutualistic and antagonistic relationships
A metadata analysis has found that habitat fragmentation greatly affects mutualistic relationships while affecting antagonistic relationships, such as predation and herbivory, to a less degree. For example, the mutualistic relationship between Mesogyne insignis and Megachile. A study has found greater pollination and increased fruit production of M. insignis in unfragmented forests verses fragmented forests. As for an example of an antagonistic relationship of nest predation, a study found that there is no increase in nest predation on fragmented forests - thus not supporting the edge effect hypothesis.
Effect on ecosystem services
Habitat fragmentation has profound effects on ecosystem services, impacting nutrient retention, species richness, and local biophysical conditions. Fragmentation-mediated processes cause generalizable responses at the population, community, and ecosystem levels, resulting in decreased nutrient retention. Furthermore, habitat fragmentation alters relationships between biodiversity and ecosystem functioning across multiple scales, affecting both the local loss of biodiversity and the local loss of function. Moreover, fragmentation can change the microclimate at both local and regional scales, influencing biodiversity through interactions with anthropogenic climate change. Overall, habitat fragmentation significantly disrupts ecosystem services by altering nutrient retention, biodiversity, and ecosystem functioning at various spatial and temporal scales.
Forest fragmentation
Forest fragmentation is a form of habitat fragmentation where forests are reduced (either naturally or man-made) to relatively small, isolated patches of forest known as forest fragments or forest remnants. The intervening matrix that separates the remaining woodland patches can be natural open areas, farmland, or developed areas. Following the principles of island biogeography, remnant woodlands act like islands of forest in a sea of pastures, fields, subdivisions, shopping malls, etc. These fragments will then begin to undergo the process of ecosystem decay.
Forest fragmentation also includes less subtle forms of discontinuities such as utility right-of-ways (ROWs). Utility ROWs are of ecological interest because they have become pervasive in many forest communities, spanning areas as large as 5 million acres in the United States. Utility ROWs include electricity transmission ROWs, gas pipeline and telecommunication ROWs. Electricity transmission ROWs are created to prevent vegetation interference with transmission lines. Some studies have shown that electricity transmission ROWs harbor more plant species than adjoining forest areas, due to alterations in the microclimate in and around the corridor. Discontinuities in forest areas associated with utility right-of-ways can serve as biodiversity havens for native bees and grassland species, as the right-of-ways are preserved in an early successional stage.
Forest fragmentation reduces food resources and habitat sources for animals thus splitting these species apart. Thus, making these animals become much more susceptible to effects of predation and making them less likely to perform interbreeding - lowering genetic diversity.
Additionally, forest fragmentation affects the native plant species present within the area by dividing large populations into smaller ones. In turn, smaller populations are more inclined to be affected by genetic drift and population performance, as well as experience increases in inbreeding activities. Moreover, fragmentation can affect the relationship present between animals and plants, such as the relationships regarding seed-dispersal or pollinator-plant relationship.
Implications
Forest fragmentation is one of the greatest threats to biodiversity in forests, especially in the tropics. The problem of habitat destruction that caused the fragmentation in the first place is compounded by:
the inability of individual forest fragments to support viable populations, especially of large vertebrates
the local extinction of species that do not have at least one fragment capable of supporting a viable population
edge effects that alter the conditions of the outer areas of the fragment, greatly reducing the amount of true forest interior habitat.
The effect of fragmentation on the flora and fauna of a forest patch depends on a) the size of the patch, and b) its degree of isolation. Isolation depends on the distance to the nearest similar patch, and the contrast with the surrounding areas. For example, if a cleared area is reforested or allowed to regenerate, the increasing structural diversity of the vegetation will lessen the isolation of the forest fragments. However, when formerly forested lands are converted permanently to pastures, agricultural fields, or human-inhabited developed areas, the remaining forest fragments, and the biota within them, are often highly isolated.
Forest patches that are smaller or more isolated will lose species faster than those that are larger or less isolated. A large number of small forest "islands" typically cannot support the same biodiversity that a single contiguous forest would hold, even if their combined area is much greater than the single forest. However, forest islands in rural landscapes greatly increase their biodiversity. In the Maulino forest of Chile fragmentation appear to not affect overall plant diversity much, and tree diversity is indeed higher in fragments than in large continuous forests.
McGill University in Montreal, Quebec, Canada released a university based newspaper statement stating that 70% of the world's remaining forest stands within one kilometre of a forest edge putting biodiversity at an immense risk based on research conducted by international scientists.
Reduced fragment area, increased isolation, and increased edge initiate changes that percolate through all ecosystems. Habitat fragmentation is able to formulate persistent outcomes which can also become unexpected such as an abundance of some species and the pattern that long temporal scales are required to discern many strong system responses.
Sustainable forest management
The presence of forest fragments influences the supply of various ecosystems in adjacent agricultural fields (Mitchell et al. 2014). Mitchell et al. (2014), researched on six varying ecosystem factors such as crop production, decomposition, pesticide regulation, carbon storage, soil fertility, and water quality regulation in soybean fields through separate distances by nearby forest fragments which all varied in isolation and size across an agricultural landscape in Quebec, Canada. Sustainable forest management can be achieved in several ways including by managing forests for ecosystem services (beyond simple provisioning), through government compensation schemes, and through effective regulation and legal frameworks. The only realistic method of conserving forests is to apply and practice sustainable forest management to risk further loss.
There is a high industrial demand for wood, pulp, paper, and other resources which the forest can provide with, thus businesses which will want more access to the cutting of forests to gain those resources. The rainforest alliance has efficiently been able to put into place an approach to sustainable forest management, and they established this in the late 1980s. Their conservation was deemed successful as it has saved over nearly half a billion acres of land around the world.
A few approaches and measures which can be taken in order to conserve forests are methods by which erosion can be minimized, waste is properly disposed, conserve native tree species to maintain genetic diversity, and setting aside forestland (provides habitat for critical wildlife species). Additionally, forest fires can also occur frequently and measures can also be taken to further prevent forest fires from occurring. For example, in Guatemala’s culturally and ecologically significant Petén region, researchers were able to find over a 20-year period, actively managed FSC-certified forests experienced substantially lower rates of deforestation than nearby protected areas, and forest fires only affected 0.1 percent of certified land area, compared to 10.4 percent of protected areas. However, it must be duly noted that short term decisions regarding forest sector employment and harvest practices can have long-term effects on biodiversity. Planted forests become increasingly important as they supply approximately a quarter of global industrial roundwood production and are predicted to account for 50% of global output within two decades (Brown, 1998; Jaakko Poyry, 1999). Although there have been many difficulties, the implementation of forest certification has been quite prominent in being able to raise effective awareness and disseminating knowledge on a holistic concept, embracing economic, environmental and social issues, worldwide. While also providing a tool for a range of other applications than assessment of sustainability, such as e.g. verifying carbon sinks.
Approaches to understanding habitat fragmentation
Two approaches are typically used to understand habitat fragmentation and its ecological impacts.
Species-oriented approach
The species-oriented approach focuses specifically on individual species and how they each respond to their environment and habitat changes with in it. This approach can be limited because it does only focus on individual species and does not allow for a broad view of the impacts of habitat fragmentation across species.
Pattern-oriented approach
The pattern-oriented approach is based on land cover and its patterning in correlation with species occurrences. One model of study for landscape patterning is the patch-matrix-corridor model developed by Richard Forman The pattern-oriented approach focuses on land cover defined by human means and activities. This model has stemmed from island biogeography and tries to infer causal relationships between the defined landscapes and the occurrence of species or groups of species within them. The approach has limitations in its collective assumptions across species or landscapes which may not account for variations amongst them.
Variegation model
The other model is the variegation model. Variegated landscapes retain much of their natural vegetation but are intermixed with gradients of modified habitat This model of habitat fragmentation typically applies to landscapes that are modified by agriculture. In contrast to the fragmentation model that is denoted by isolated patches of habitat surrounded by unsuitable landscape environments, the variegation model applies to landscapes modified by agriculture where small patches of habitat remain near the remnant original habitat. In between these patches are a matrix of grassland that is often modified versions of the original habitat. These areas do not present as much of a barrier to native species.
| Biology and health sciences | Ecology | Biology |
1260189 | https://en.wikipedia.org/wiki/Arizona%20bark%20scorpion | Arizona bark scorpion | The Arizona bark scorpion (Centruroides sculpturatus, once included in Centruroides exilicauda) is a small light brown scorpion common to the Sonoran Desert in the southwestern United States and northwestern Mexico. An adult male can reach of body length, while a female is slightly smaller, with a maximum length of .
Predators
Arizona bark scorpions are eaten by a wide variety of animals such as pallid bats, birds (especially owls), reptiles, and other vertebrates. Some examples include spiders, snakes, peccaries, rodents, and other scorpions. Development, pesticides and collecting scorpions for research or the pet trade also reduces the bark scorpion population.
The painful and potentially deadly venom of Arizona bark scorpions has little effect on grasshopper mice. Scientists have found the scorpion toxin acts as an analgesic rather than a pain stimulant in grasshopper mice.
Life cycle
Arizona bark scorpions have a gestation period of several months, are born live, and are gently guided onto their mother's back. The female usually gives birth to anywhere from 25 to 35 young. These remain with their mother until their first molt, which can be up to three weeks after birth. Arizona bark scorpions have a life expectancy of about six years.
Arizona bark scorpions, like most other scorpions, are incredibly resilient. During US nuclear testing, scorpions, along with cockroaches and lizards, were found near ground zero with no recorded adverse effects.
Habitat
The Arizona bark scorpion is nocturnal, and particularly well adapted to the desert: layers of wax on its exoskeleton make it resistant to water loss. Nevertheless, Arizona bark scorpions hide during the heat of the day, typically under rocks, wood piles, or tree bark. Arizona bark scorpions do burrow, and are commonly found in homes, requiring only 1/16 of an inch for entry.
Arizona bark scorpions prefer riparian areas with mesquite, cottonwood, and sycamore groves, all of which have sufficient moisture and humidity to support insects and other prey species. The popularity of irrigated lawns, and other systems which increase environmental humidity in residential areas, has led to a massive increase in the number of these animals in some areas.
Centruroides scorpions are unusual in that they are the only genus in the Southwest that can climb walls, trees, and other objects with a sufficiently rough surface. Arizona bark scorpions practice negative geotaxis, preferring an upside down orientation, which often results in people being stung due to the scorpion being on the underside of an object.
The Arizona bark scorpion preys on small and medium-sized animals such as beetles, spiders, crickets, cockroaches, other insects and other scorpions.
The range of the Arizona bark scorpion is from southern California, southern Arizona, southern Nevada, extreme southwestern Utah and western New Mexico. They are also found in Baja California, Sonora and Chihuahua, Mexico.
Venom
The Arizona bark scorpion is the most venomous scorpion in North America, and its venom can cause severe pain (coupled with numbness, tingling, and vomiting) in adult humans, typically lasting between 24 and 72 hours. Temporary dysfunction in the area stung is common; e.g. a hand or possibly arm can be immobilized or experience convulsions. It also may cause loss of breath for a short time. Due to the extreme pain induced, many victims describe sensations of electrical jolts after envenomation. Two recorded fatalities have occurred in the state of Arizona since 1968; the number of victims stung each year in Arizona and New Mexico is estimated to be in the thousands.
Antivenom
An antivenom was developed for this species at Arizona State University by Dr. Herbert L. Stahnke, and produced in quantities sufficient to treat individuals within the state of Arizona. This antivenom was not FDA approved, but use within the state of Arizona was allowable and very successful in shortening the duration of symptoms and hospitalization. Production of this antivenom ceased by 2000 and the product was unavailable by 2004. A Mexican-produced antivenom, Anascorp [Antivenin Centruroides (scorpion) F(ab′)2, Laboratorios Silanes, Instituto Bioclon SA de CV], received FDA approval on August 3, 2011, and is now in use.
First aid
Basic first aid measures can be used to help mediate Arizona bark scorpion stings:
Clean sting site with soap and water
Apply a cool compress (cool cloth)
Take acetaminophen (paracetamol) or ibuprofen for local pain and swelling
Medical emergencies
Arizona poison control centers suggest immediate medical attention if severe symptoms occur, particularly in young children. The Poison Center may be reached at 1-800-222-1222.
UV lighting
Arizona bark scorpions, like most other scorpions, will glow when exposed to a blacklight. This is particularly useful in scorpion detection, since Arizona bark scorpions are active during the night, and can be easily spotted using this method. Typical UV LED flashlights enable their human operator to readily detect Arizona bark scorpions at a distance of approximately six feet. Newly molted Arizona bark scorpions will not glow under ultraviolet light for a few days after molting.
Control and prevention
Arizona bark scorpions are tan or light beige tone in color and very small, making them difficult to detect especially on natural terrain (rocky land, multiple vegetation and soil textured land). They often look for places to hide, and they will not seek out humans unless provoked or defending young. Several methods of control have historically been used to control Arizona bark scorpions, such as physical barriers (scorpions are unable to climb smooth surfaces), pesticides, glue boards, and removing any scorpion congregation areas in the vicinity of the building.
| Biology and health sciences | Scorpions | Animals |
1261615 | https://en.wikipedia.org/wiki/Rectangular%20function | Rectangular function | The rectangular function (also known as the rectangle function, rect function, Pi function, Heaviside Pi function, gate function, unit pulse, or the normalized boxcar function) is defined as
Alternative definitions of the function define to be 0, 1, or undefined.
Its periodic version is called a rectangular wave.
History
The rect function has been introduced by Woodward in as an ideal cutout operator, together with the sinc function as an ideal interpolation operator, and their counter operations which are sampling (comb operator) and replicating (rep operator), respectively.
Relation to the boxcar function
The rectangular function is a special case of the more general boxcar function:
where is the Heaviside step function; the function is centered at and has duration , from to
Fourier transform of the rectangular function
The unitary Fourier transforms of the rectangular function are
using ordinary frequency , where is the normalized form of the sinc function and
using angular frequency , where is the unnormalized form of the sinc function.
For , its Fourier transform isNote that as long as the definition of the pulse function is only motivated by its behavior in the time-domain experience, there is no reason to believe that the oscillatory interpretation (i.e. the Fourier transform function) should be intuitive, or directly understood by humans. However, some aspects of the theoretical result may be understood intuitively, as finiteness in time domain corresponds to an infinite frequency response. (Vice versa, a finite Fourier transform will correspond to infinite time domain response.)
Relation to the triangular function
We can define the triangular function as the convolution of two rectangular functions:
Use in probability
Viewing the rectangular function as a probability density function, it is a special case of the continuous uniform distribution with The characteristic function is
and its moment-generating function is
where is the hyperbolic sine function.
Rational approximation
The pulse function may also be expressed as a limit of a rational function:
Demonstration of validity
First, we consider the case where Notice that the term is always positive for integer However, and hence approaches zero for large
It follows that:
Second, we consider the case where Notice that the term is always positive for integer However, and hence grows very large for large
It follows that:
Third, we consider the case where We may simply substitute in our equation:
We see that it satisfies the definition of the pulse function. Therefore,
Dirac delta function
The rectangle function can be used to represent the Dirac delta function . Specifically,For a function , its average over the width around 0 in the function domain is calculated as,
To obtain , the following limit is applied,
and this can be written in terms of the Dirac delta function as,
The Fourier transform of the Dirac delta function is
where the sinc function here is the normalized sinc function. Because the first zero of the sinc function is at and goes to infinity, the Fourier transform of is
means that the frequency spectrum of the Dirac delta function is infinitely broad. As a pulse is shorten in time, it is larger in spectrum.
| Mathematics | Specific functions | null |
1261767 | https://en.wikipedia.org/wiki/Chlorophytum%20comosum | Chlorophytum comosum | Chlorophytum comosum, usually called spider plant or common spider plant due to its spider-like look, also known as spider ivy, airplane plant, ribbon plant (a name it shares with Dracaena sanderiana), and hen and chickens, is a species of evergreen perennial flowering plant of the family Asparagaceae. It is native to tropical and Southern Africa but has become naturalized in other parts of the world, including Western Australia and Bangladesh. Chlorophytum comosum is easy to grow as a houseplant because of its resilience, but it can be sensitive to the fluoride in tap water, which commonly gives it "burnt tips". Variegated forms are the most popular.
Description
Chlorophytum comosum grows to about tall, although as a hanging plant it can descend many feet. It has fleshy, tuberous roots, each about long. The long narrow leaves reach a length of and are around wide.
Flowers are produced in a long, branched inflorescence, which can reach a length of up to and eventually bends downward to meet the earth. Flowers initially occur in clusters of 1–6 at intervals along the stem (scape) of the inflorescence. Each cluster is at the base of a bract, which ranges from in length, becoming smaller toward the end of the inflorescence. Most of the flowers that are produced initially die off, so that relatively, the inflorescences are sparsely flowered.
Individual flowers are greenish-white, borne on stalks (pedicels) some long. Each flower has six triply veined tepals that are long and slightly hooded or boat-shaped at their tips. The stamens consist of a pollen-producing anther about long with a filament of similar length or slightly longer. The central style is long. Seeds are produced in a capsule, long, on stalks (pedicels) that lengthen to up to .
The inflorescences carry not only flowers but also vegetative plantlets at the tips of their branches, which eventually droop and touch the soil, developing adventitious roots. The stems (scapes) of the inflorescence are called "stolons" in some sources, but this term is more correctly used for stems that do not bear flowers and have roots at the nodes.
Taxonomy
The first formal description of Chlorophytum comosum was by the Swedish naturalist Carl Peter Thunberg as Anthericum comosum in the 1794 volume of Prodromus Plantarum Capensium, Thunberg's work on the plants of South Africa. The species was subsequently moved to a number of different genera, including Phalangium, Caesia, Hartwegia Nees, and Hollia, before receiving its current placement in Chlorophytum by Jacques in 1862.
The species has been confused with Chlorophytum capense (L.) Voss by some authors, but this is a different species.
Intra-specific variation
There are three described varieties of the species: the autonym C. comosum var. comosum has strap-shaped narrow leaves and is found along forest margins; C. comosum var. bipindense has broader, petiolate leaves with stripes on the underside and the inflorescences are 2–3 times the length of the leaves; and C. comosum var. sparsiflorum also has broader leaves that narrow to the base, and usually lacks a petiole and the striping on the underside of the leaf and the inflorescences are up to two times the length of the leaves. The latter two are rainforest-dwelling taxa that had been described earlier as separate species, but botanists Axel Dalberg Poulsen and Inger Nordal reduced the taxa to varieties of C. comosum in 2005.
Delimitation of species boundaries within the genus Chlorophytum is reported to be difficult, possibly because of several evolutionary radiations into forest environments that led to morphological aspects that are too similar to reliably distinguish separate species. The evidence given to support this is the widespread distribution of most taxa in the genus and poor seed dispersal, leading to the conclusion of deeper evolutionary divergence among the taxa.
The three described varieties in C. comosum could be an example of this convergent evolution of leaf shape among the forest-dwelling varieties from species of disparate origin, leading to the species C. comosum being polyphyletic, instead of the traditional view of morphological divergence among the varieties within the species with the assumption of a common origin (monophyly). The widespread C. comosum var. comosum has slender, near linear leaves that lack a petiole similar to plants found in cultivation and is only found growing at the margins of the rainforest. The two other varieties, C. comosum var. sparsiflorum and C. comosum var. bipindense, possess petioles and have broader leaves necessary for collecting more light in the shady Guineo-Congolean rainforest. A study published in 2005 used 16 morphological characters and was unable to delimit species boundaries among these three taxa, so they were relegated to varietal status. A follow-up study published in 2008 provided preliminary evidence from phylogenetic analysis of plastid and nuclear DNA sequences that established samples from disparate collections sites identified as C. comosum were polyphyletic.
Distribution
Chlorophytum comosum has a widespread native distribution in Africa, being native to six of the ten World Geographical Scheme for Recording Plant Distributions regions of Africa (West Tropical Africa, West-Central Tropical Africa, Northeast Tropical Africa, East Tropical Africa, South Tropical Africa, and Southern Africa).
Cultivation
Chlorophytum comosum is a popular houseplant. The species with all-green leaves forms only a small proportion of plants sold. More common are two variegated cultivars:
C. comosum 'Vittatum' has mid-green leaves with a broad central white stripe. It is often sold in hanging baskets to display the plantlets. The long stems are white. There is also a "curly" version with this type of striping and compact size.
C. comosum 'Variegatum' has darker green leaves with white margins. It is generally smaller than the previous cultivar. The long stems are green.
Both cultivars have gained the Royal Horticultural Society's Award of Garden Merit (confirmed 2017).
In 2021, 17 other cultivars were listed, including 'Bonnie', 'Green Bonnie' and 'Hawaiian'.
Propagation
Propagating Chlorophytum comosum commonly occurs through potting the plantlets, informally referred to as 'spiderettes', or 'pups', directly into potting soil, or pumice, attached to the main plant or cutting the running stems and then potting them.
Spider plants are easy to grow, being able to thrive in a wide range of conditions. They will tolerate temperatures down to , but grow best at temperatures between and . Plants can be damaged by high fluoride or boron levels.
Toxicity and effects on pets
Spider plants are non-toxic to humans and pets, and are considered edible.
Air purification
The NASA Clean Air Study suggested that air plants were effective at removing common household air toxins formaldehyde and xylene; however, these results are not applicable to typical buildings, where outdoor-to-indoor air exchange already removes volatile organic compounds (VOCs) at a rate that could only be matched by the placement of 10–1000 plants/m of a building's floor space.
The results also failed to replicate in future studies, with a 2014 review stating that:
In the laboratory settings uses in the Clean Air Study, spider plants were shown to reduce formaldehyde pollution, and approximately 70 plants would neutralize the formaldehyde released by materials in a representative (c. ) energy-efficient house, assuming each plant occupies a pot.
Edibility
The tuberous roots are reportedly edible (whether raw or cooked unstated) although mild laxative effects are claimed by the Nguni people of its native South Africa.
| Biology and health sciences | Asparagales | Plants |
1262082 | https://en.wikipedia.org/wiki/Cistern | Cistern | A cistern (; , ; ) is a waterproof receptacle for holding liquids, usually water. Cisterns are often built to catch and store rainwater. To prevent leakage, the interior of the cistern is often lined with hydraulic plaster.
Cisterns are distinguished from wells by their waterproof linings. Modern cisterns range in capacity from a few litres to thousands of cubic meters, effectively forming covered reservoirs.
Origins
Early domestic and agricultural use
Waterproof lime plaster cisterns in the floors of houses are features of Neolithic village sites of the Levant at, for instance, Ramad and Lebwe, and by the late fourth millennium BC, as at Jawa in northeastern Lebanon, cisterns are essential elements of emerging water management techniques in dry-land farming communities.
Early examples of ancient cisterns, found in Israel, include a significant discovery at Tel Hazor, where a large cistern was carved into bedrock beneath a palace dating to the Late Bronze Age. Similar systems were uncovered at Ta'anakh. In the Iron Age, underground water systems were constructed in royal centers and settlements throughout ancient Israel, marking some of the earliest instances of engineering activity in urban planning.
The Ancient Roman impluvium, a standard feature of the domus house, generally had a cistern underneath. The impluvium and associated structures collected, filtered, cooled, and stored the water, and also cooled and ventilated the house.
Castle cisterns
In the Middle Ages, cisterns were often constructed in hill castles in Europe, especially where wells could not be dug deeply enough. There were two types: the tank cistern and the filter cistern. Such a filter cistern was built at the Riegersburg in Austrian Styria, where a cistern was hewn out of the lava rock. Rain water passed through a sand filter and collected in the cistern. The filter cleaned the rain water and enriched it with minerals.
Present-day use
Cisterns are commonly prevalent in areas where water is scarce, either because it is rare or has been depleted due to heavy use. Historically, the water was used for many purposes including cooking, irrigation, and washing. Present-day cisterns are often used only for irrigation due to concerns over water quality. Cisterns today can also be outfitted with filters or other water purification methods when the water is intended for consumption. It is not uncommon for a cistern to be open in some manner in order to catch rain or to include more elaborate rainwater harvesting systems. It is important in these cases to have a system that does not leave the water open to algae or to mosquitoes, which are attracted to the water and then potentially carry disease to nearby humans.
One particularly unique modern utilization of cisterns is found in San Francisco, which has historically been subject to devastating fires. As a precautionary measure, in 1850, funds were allocated to construct over 100 cisterns across the city to be utilized in case of fire. The city's firefighting network, the Auxiliary Water Supply System (AWSS) maintains a network of 177 independent underground water cisterns, with sizes varying from 75,000 US gallons (280,000 L) to over 200,000 US gallons (760,000 L) depending on location with a total storage capacity of over 11 million U.S. gallons (42 million liters) of water. These cisterns are easily spotted at street level with manholes labeled CISTERN S.F.F.D surrounded by red brick circles or rectangles. The cisterns are completely separate from the rest of the city’s water supply, ensuring that in the event of an earthquake, additional backup is available regardless of the condition of the city's mainline water system.
Some cisterns sit on the top of houses or on the ground higher than the house, and supply the running water needs for the house. They are often supplied by wells with electric pumps, or are filled manually or by truck delivery, rather than by rainwater collection. Very common throughout Brazil, for example, they were traditionally made of concrete walls (much like the houses themselves), with a similar concrete top (about 5 cm/2 inches thick), with a piece that can be removed for water filling and then reinserted to keep out debris and insects. Modern cisterns are manufactured out of plastic (in Brazil with a characteristic bright blue color, round, in capacities of about 10,000 and 50,000 liters (2641 and 13,208 gallons)). These cisterns differ from water tanks in the sense that they are not entirely enclosed and sealed with one form, rather they have a lid made of the same material as the cistern, which is removable by the user.
To keep a clean water supply, the cistern must be kept clean. It is important to inspect them regularly, keep them well enclosed, and to occasionally empty and clean them with a proper dilution of chlorine and to rinse them well. Well water must be inspected for contaminants coming from the ground source. City water has up to 1ppm (parts per million) chlorine added to the water to keep it clean. If there is any question about the water supply at any point (source to tap), then the cistern water should not be used for drinking or cooking. If it is of acceptable quality and consistency, then it can be used for (1) toilets, and housecleaning; (2) showers and handwashing; (3) washing dishes, with proper sanitation methods, and for the highest quality, (4) cooking and drinking. Water of non-acceptable quality for the aforementioned uses may still be used for irrigation. If it is free of particulates but not low enough in bacteria, then boiling may also be an effective method to prepare the water for drinking.
Many greenhouses rely on a cistern to help meet their water needs, particularly in the United States. Some countries or regions, such as Flanders, Bermuda and the U.S. Virgin Islands, have strict laws requiring that rainwater harvesting systems be built alongside any new construction, and cisterns can be used in these cases. In Bermuda, for example, its familiar white-stepped roofs seen on houses are part of the rainwater collection system, where water is channeled by roof gutters to below-ground cisterns. Other countries, such as Japan, Germany, and Spain, also offer financial incentives or tax credit for installing cisterns. Cisterns may also be used to store water for firefighting in areas where there is an inadequate water supply. The city of San Francisco, notably, maintains fire cisterns under its streets in case the primary water supply is disrupted. In many flat areas, the use of cisterns is encouraged to absorb excess rainwater which otherwise can overload sewage or drainage systems by heavy rains (certainly in urban areas where a lot of ground is surfaced and doesn't let the ground absorb water).
Bathing
In some southeast Asian countries such as Malaysia and Indonesia showers are traditionally taken by pouring water over one's body with a dipper (this practice comes from before piped water was common). Many bathrooms even in modern houses are constructed with a small cistern to hold water for bathing by this method.
Toilet cisterns
The modern toilet utilises a cistern to reserve and hold the correct amount of water required to flush the toilet bowl. In earlier toilets, the cistern was located high above the toilet bowl and connected to it by a long pipe. It was necessary to pull a hanging chain connected to a release valve located inside the cistern in order to flush the toilet. Modern toilets may be close coupled, with the cistern mounted directly on the toilet bowl and no intermediate pipe. In this arrangement, the flush mechanism (lever or push button) is usually mounted on the cistern. Concealed cistern toilets, where the cistern is built into the wall behind the toilet, are also available. A flushing trough is a type of cistern used to serve more than one WC pan at one time. These cisterns are becoming less common, however. The cistern was the genesis of the modern bidet.
At the beginning of the flush cycle, as the water level in the toilet cistern tank drops, the flush valve flapper falls back to the bottom, stopping the main flow to the flush tube. Because the tank water level has yet to reach the fill line, water continues to flow from the tank and bowl fill tubes. When the water again reaches the fill line, the float will release the fill valve shaft and water flow will stop.
One Million Cisterns Program
In Northeastern Brazil, the One Million Cisterns Program (Programa 1 Milhão de Cisternas or P1MC) has assisted local people with water management. The Brazilian government adopted this new policy of rainwater harvesting in 2013. The Semi-Arid Articulation (ASA) has been providing managerial and technological support to establish cement-layered containers, called cisterns, to harvest and store rainwater for small farm-holders in 34 territories of nine states where ASA operates (Minas Gerais, Bahia, Sergipe, Alagoas, Pernambuco, Paraíba, Rio Grande do Norte, Ceará and Piauí).
The rainwater falling on the rooftops is directed through pipelines or gutters and stored in the cistern. The cistern is covered with a lid to avoid evaporation. Each cistern has a capacity of 16,000 liters. Water collected in it during 3–4 months of the rainy season can sustain the requirement for drinking, cooking, and other basic sanitation purposes for rest of the dry periods. By 2016, 1.2 million rainwater harvesting cisterns were implemented for human consumption alone. After positive results of P1MC, the government introduced another program named "One Land, Two Water Program" (Uma Terra, Duas Águas, P1 + 2), which provides a farmer with another slab cistern to support agricultural production.
Notable examples
Basilica Cistern in Istanbul, Turkey
Aljibe of the in Cáceres, Spain
Portuguese cistern (Mazagan) in El Jadida, Morocco
Cistern in Silves, Portugal
Matera, southern Italy
Asa of Judah built a cistern. The prophet Jeremiah was later thrown in it after prophesying the Babylonian invasion
Cistern in Genesis 37:20, 22
| Technology | Hydraulic infrastructure | null |
1262556 | https://en.wikipedia.org/wiki/Methylamine | Methylamine | Methylamine, also known as methanamine, is an organic compound with a formula of . This colorless gas is a derivative of ammonia, but with one hydrogen atom being replaced by a methyl group. It is the simplest primary amine.
Methylamine is sold as a solution in methanol, ethanol, tetrahydrofuran, or water, or as the anhydrous gas in pressurized metal containers. Industrially, methylamine is transported in its anhydrous form in pressurized railcars and tank trailers. It has a strong odor similar to rotten fish. Methylamine is used as a building block for the synthesis of numerous other commercially available compounds.
Industrial production
Methylamine has been produced industrially since the 1920s (originally by Commercial Solvents Corporation for dehairing of animal skins). This was made possible by and his wife Eugenia who discovered amination of alcohols, including methanol, on alumina or kaolin catalyst after WWI, filed two patent applications in 1919 and published an article in 1921.
It is now prepared commercially by the reaction of ammonia with methanol in the presence of an aluminosilicate catalyst. Dimethylamine and trimethylamine are co-produced; the reaction kinetics and reactant ratios determine the ratio of the three products. The product most favored by the reaction kinetics is trimethylamine.
In this way, an estimated 115,000 tons were produced in 2005.
Laboratory methods
Methylamine was first prepared in 1849 by Charles-Adolphe Wurtz via the hydrolysis of methyl isocyanate and related compounds. An example of this process includes the use of the Hofmann rearrangement, to yield methylamine from acetamide and bromine.
In the laboratory, methylamine hydrochloride is readily prepared by various other methods. One method entails treating formaldehyde with ammonium chloride.
The colorless hydrochloride salt can be converted to an amine by the addition of a strong base, such as sodium hydroxide (NaOH):
Another method entails reducing nitromethane with zinc and hydrochloric acid.
Another method of methylamine production is spontaneous decarboxylation of glycine with a strong base in water.
Reactivity and applications
Methylamine is a good nucleophile as it is an unhindered amine. As an amine it is considered a weak base. Its use in organic chemistry is pervasive. Some reactions involving simple reagents include: with phosgene to methyl isocyanate, with carbon disulfide and sodium hydroxide to the sodium methyldithiocarbamate, with chloroform and base to methyl isocyanide and with ethylene oxide to methylethanolamines. Liquid methylamine has solvent properties analogous to those of liquid ammonia.
Representative commercially significant chemicals produced from methylamine include the pharmaceuticals ephedrine and theophylline, the pesticides carbofuran, carbaryl, and metham sodium, and the solvents N-methylformamide and N-methylpyrrolidone. The preparation of some surfactants and photographic developers require methylamine as a building block.
Biological chemistry
Methylamine arises as a result of putrefaction and is a substrate for methanogenesis.
Additionally, methylamine is produced during PADI4-dependent arginine demethylation.
Safety
The LD50 (mouse, s.c.) is 2.5 g/kg.
The Occupational Safety and Health Administration (OSHA) and National Institute for Occupational Safety and Health (NIOSH) have set occupational exposure limits at 10 ppm or 12 mg/m3 over an eight-hour time-weighted average.
Regulation
In the United States, methylamine is controlled as a List 1 precursor chemical by the Drug Enforcement Administration due to its use in the illicit production of methamphetamine.
In popular culture
Fictional characters Walter White and Jesse Pinkman use aqueous methylamine as part of a process to synthesize methamphetamine in the AMC series Breaking Bad.
| Physical sciences | Hydrogen compounds | Chemistry |
23694864 | https://en.wikipedia.org/wiki/Gamma-ray%20astronomy | Gamma-ray astronomy | Gamma-ray astronomy is a subfield of astronomy where scientists observe and study celestial objects and phenomena in outer space which emit cosmic electromagnetic radiation in the form of gamma rays, i.e. photons with the highest energies (above 100 keV) at the very shortest wavelengths. Radiation below 100 keV is classified as X-rays and is the subject of X-ray astronomy.
In most cases, gamma rays from solar flares and Earth's atmosphere fall in the MeV range, but it's now known that solar flares can also produce gamma rays in the GeV range, contrary to previous beliefs. Much of the detected gamma radiation stems from collisions between hydrogen gas and cosmic rays within our galaxy. These gamma rays, originating from diverse mechanisms such as electron-positron annihilation, the inverse Compton effect and in some cases gamma decay, occur in regions of extreme temperature, density, and magnetic fields, reflecting violent astrophysical processes like the decay of neutral pions. They provide insights into extreme events like supernovae, hypernovae, and the behavior of matter in environments such as pulsars and blazars. A huge number of gamma ray emitting high-energy systems like black holes, stellar coronas, neutron stars, white dwarf stars, remnants of supernova, clusters of galaxies, including the Crab Nebula and the Vela pulsar (the most powerful source so far), have been identified, alongside an overall diffuse gamma-ray background along the plane of the Milky Way galaxy. Cosmic radiation with the highest energy triggers electron-photon cascades in the atmosphere, while lower-energy gamma rays are only detectable above it. Gamma-ray bursts, like GRB 190114C, are transient phenomena challenging our understanding of high-energy astrophysical processes, ranging from microseconds to several hundred seconds.
Gamma rays are difficult to detect due to their high energy and their blocking by the Earth’s atmosphere, necessitating balloon-borne detectors and artificial satellites in space. Early experiments in the 1950s and 1960s used balloons to carry instruments to access altitudes where the atmospheric absorption of gamma rays is low, followed by the launch of the first gamma-ray satellites: SAS 2 (1972) and COS-B (1975). These were defense satellites originally designed to detect gamma rays from secret nuclear testing, but they luckily discovered puzzling gamma-ray bursts coming from deep space. In the 1970s, satellite observatories found several gamma-ray sources, among which a very strong source called Geminga was later identified as a pulsar in proximity. The Compton Gamma Ray Observatory (launched in 1991) revealed numerous gamma-ray sources in space. Today, both ground-based observatories like the VERITAS array and space-based telescopes like the Fermi Gamma-ray Space Telescope (launched in 2008) contribute significantly to gamma-ray astronomy. This interdisciplinary field involves collaboration among physicists, astrophysicists, and engineers in projects like the High Energy Stereoscopic System (H.E.S.S.), which explores extreme astrophysical environments like the vicinity of black holes in active galactic nuclei.
Studying gamma rays provides valuable insights into extreme astrophysical environments, as observed by the H.E.S.S. Observatory. Ongoing research aims to expand our understanding of gamma-ray sources, such as blazars, and their implications for cosmology. As GeV gamma rays are important in the study of extra-solar, and especially extragalactic, astronomy, new observations may complicate some prior models and findings.
Future developments in gamma-ray astronomy will integrate data from gravitational wave and neutrino observatories (Multi-messenger astronomy), enriching our understanding of cosmic events like neutron star mergers. Technological advancements, including advanced mirror designs, better camera technologies, improved trigger systems, faster readout electronics, high-performance photon detectors like Silicon photomultipliers (SiPMs), alongside innovative data processing algorithms like time-tagging techniques and event reconstruction methods, will enhance spatial and temporal resolution. Machine learning algorithms and big data analytics will facilitate the extraction of meaningful insights from vast datasets, leading to discoveries of new gamma-ray sources, identification of specific gamma-ray signatures, and improved modeling of gamma-ray emission mechanisms. Future missions may include space telescopes and lunar gamma-ray observatories (taking advantage of the Moon's lack of atmosphere and stable environment for prolonged observations), enabling observations in previously inaccessible regions. The ground-based Cherenkov Telescope Array project, a next-generation gamma ray observatory which will incorporate many of these improvements and will be ten times more sensitive, is planned to be fully operational by 2025.
Early history
Long before experiments could detect gamma rays emitted by cosmic sources, scientists had known that the universe should be producing them. Work by Eugene Feenberg and Henry Primakoff in 1948, Sachio Hayakawa and I.B. Hutchinson in 1952, and, especially, Philip Morrison in 1958 had led scientists to believe that a number of different processes which were occurring in the universe would result in gamma-ray emission. These processes included cosmic ray interactions with interstellar gas, supernova explosions, and interactions of energetic electrons with magnetic fields.
However, it was not until the 1960s that our ability to actually detect these emissions came to pass.
Most gamma rays coming from space are absorbed by the Earth's atmosphere, so gamma-ray astronomy could not develop until it was possible to get detectors above all or most of the atmosphere using balloons and spacecraft. The first gamma-ray telescope carried into orbit, on the Explorer 11 satellite in 1961, picked up fewer than 100 cosmic gamma-ray photons. They appeared to come from all directions in the Universe, implying some sort of uniform "gamma-ray background". Such a background would be expected from the interaction of cosmic rays (very energetic charged particles in space) with interstellar gas.
The first true astrophysical gamma-ray sources were solar flares, which revealed the strong 2.223 MeV line predicted by Morrison. This line results from the formation of deuterium via the union of a neutron and proton; in a solar flare the neutrons appear as secondaries from interactions of high-energy ions accelerated in the flare process. These first gamma-ray line observations were from OSO 3, OSO 7, and the Solar Maximum Mission, the latter spacecraft launched in 1980. The solar observations inspired theoretical work by Reuven Ramaty and others.
Significant gamma-ray emission from our galaxy was first detected in 1967 by the detector aboard the OSO 3 satellite. It detected 621 events attributable to cosmic gamma rays. However, the field of gamma-ray astronomy took great leaps forward with the SAS-2 (1972) and the Cos-B (1975–1982) satellites. These two satellites provided an exciting view into the high-energy universe (sometimes called the 'violent' universe, because the kinds of events in space that produce gamma rays tend to be high-speed collisions and similar processes). They confirmed the earlier findings of the gamma-ray background, produced the first detailed map of the sky at gamma-ray wavelengths, and detected a number of point sources. However the resolution of the instruments was insufficient to identify most of these point sources with specific visible stars or stellar systems.
A discovery in gamma-ray astronomy came in the late 1960s and early 1970s from a constellation of military defense satellites. Detectors on board the Vela satellite series, designed to detect flashes of gamma rays from nuclear bomb blasts, began to record bursts of gamma rays from deep space rather than the vicinity of the Earth. Later detectors determined that these gamma-ray bursts are seen to last for fractions of a second to minutes, appearing suddenly from unexpected directions, flickering, and then fading after briefly dominating the gamma-ray sky. Studied since the mid-1980s with instruments on board a variety of satellites and space probes, including Soviet Venera spacecraft and the Pioneer Venus Orbiter, the sources of these enigmatic high-energy flashes remain a mystery. They appear to come from far away in the Universe, and currently the most likely theory seems to be that at least some of them come from so-called hypernova explosions—supernovas creating black holes rather than neutron stars.
Nuclear gamma rays were observed from the solar flares of August 4 and 7, 1972, and November 22, 1977.
A solar flare is an explosion in a solar atmosphere and was originally detected visually in the Sun. Solar flares create massive amounts of radiation across the full electromagnetic spectrum from the longest wavelength, radio waves, to high energy gamma rays. The correlations of the high energy electrons energized during the flare and the gamma rays are mostly caused by nuclear combinations of high energy protons and other heavier ions. These gamma rays can be observed and allow scientists to determine the major results of the energy released, which is not provided by the emissions from other wavelengths.
| Physical sciences | High-energy astronomy | Astronomy |
19721986 | https://en.wikipedia.org/wiki/Directed%20graph | Directed graph | In mathematics, and more specifically in graph theory, a directed graph (or digraph) is a graph that is made up of a set of vertices connected by directed edges, often called arcs.
Definition
In formal terms, a directed graph is an ordered pair where
V is a set whose elements are called vertices, nodes, or points;
A is a set of ordered pairs of vertices, called arcs, directed edges (sometimes simply edges with the corresponding set named E instead of A), arrows, or directed lines.
It differs from an ordinary or undirected graph, in that the latter is defined in terms of unordered pairs of vertices, which are usually called edges, links or lines.
The aforementioned definition does not allow a directed graph to have multiple arrows with the same source and target nodes, but some authors consider a broader definition that allows directed graphs to have such multiple arcs (namely, they allow the arc set to be a multiset). Sometimes these entities are called directed multigraphs (or multidigraphs).
On the other hand, the aforementioned definition allows a directed graph to have loops (that is, arcs that directly connect nodes with themselves), but some authors consider a narrower definition that does not allow directed graphs to have loops.
Directed graphs without loops may be called simple directed graphs, while directed graphs with loops may be called loop-digraphs (see section Types of directed graph).
Types of directed graphs
Subclasses
Symmetric directed graphs are directed graphs where all edges appear twice, one in each direction (that is, for every arrow that belongs to the digraph, the corresponding inverse arrow also belongs to it). (Such an edge is sometimes called "bidirected" and such graphs are sometimes called "bidirected", but this conflicts with the meaning for bidirected graphs.)
Simple directed graphs are directed graphs that have no loops (arrows that directly connect vertices to themselves) and no multiple arrows with same source and target nodes. As already introduced, in case of multiple arrows the entity is usually addressed as directed multigraph. Some authors describe digraphs with loops as loop-digraphs.
Complete directed graphs are simple directed graphs where each pair of vertices is joined by a symmetric pair of directed arcs (it is equivalent to an undirected complete graph with the edges replaced by pairs of inverse arcs). It follows that a complete digraph is symmetric.
Semicomplete multipartite digraphs are simple digraphs in which the vertex set is partitioned into sets such that for every pair of vertices x and y in different sets, there is an arc between x and y. There can be one arc between x and y or two arcs in opposite directions.
Semicomplete digraphs are simple digraphs where there is an arc between each pair of vertices. Every semicomplete digraph is a semicomplete multipartite digraph in a trivial way, with each vertex constituting a set of the partition.
Quasi-transitive digraphs are simple digraphs where for every triple x, y, z of distinct vertices with arcs from x to y and from y to z, there is an arc between x and z. There can be just one arc between x and z or two arcs in opposite directions. A semicomplete digraph is a quasi-transitive digraph. There are extensions of quasi-transitive digraphs called k-quasi-transitive digraphs.
Oriented graphs are directed graphs having no opposite pairs of directed edges (i.e. at most one of and may be arrows of the graph). It follows that a directed graph is an oriented graph if and only if it has no 2-cycle. Such a graph can be obtained by applying an orientation to an undirected graph.
Tournaments are oriented graphs obtained by choosing a direction for each edge in undirected complete graphs. A tournament is a semicomplete digraph.
A directed graph is acyclic if it has no directed cycles. The usual name for such a digraph is directed acyclic graph (DAG).
Multitrees are DAGs in which there are no two distinct directed paths from the same starting vertex to the same ending vertex.
Oriented trees or polytrees are DAGs formed by orienting the edges of trees (connected, acyclic undirected graphs).
Rooted trees are oriented trees in which all edges of the underlying undirected tree are directed either away from or towards the root (they are called, respectively, arborescences or out-trees, and in-trees.
Digraphs with supplementary properties
Weighted directed graphs (also known as directed networks) are (simple) directed graphs with weights assigned to their arrows, similarly to weighted graphs (which are also known as undirected networks or weighted networks).
Flow networks are weighted directed graphs where two nodes are distinguished, a source and a sink.
Rooted directed graphs (also known as flow graphs) are digraphs in which a vertex has been distinguished as the root.
Control-flow graphs are rooted digraphs used in computer science as a representation of the paths that might be traversed through a program during its execution.
Signal-flow graphs are directed graphs in which nodes represent system variables and branches (edges, arcs, or arrows) represent functional connections between pairs of nodes.
Flow graphs are digraphs associated with a set of linear algebraic or differential equations.
State diagrams are directed multigraphs that represent finite-state machines.
Commutative diagrams are digraphs used in category theory, where the vertices represent (mathematical) objects and the arrows represent morphisms, with the property that all directed paths with the same start and endpoints lead to the same result by composition.
In the theory of Lie groups, a quiver Q is a directed graph serving as the domain of, and thus characterizing the shape of, a representation V defined as a functor, specifically an object of the functor category FinVctKF(Q) where F(Q) is the free category on Q consisting of paths in Q and FinVctK is the category of finite-dimensional vector spaces over a field K. Representations of a quiver label its vertices with vector spaces and its edges (and hence paths) compatibly with linear transformations between them, and transform via natural transformations.
Basic terminology
An arc is considered to be directed from x to y; y is called the head and x is called the tail of the arc; y is said to be a direct successor of x and x is said to be a direct predecessor of y. If a path leads from x to y, then y is said to be a successor of x and reachable from x, and x is said to be a predecessor of y. The arc is called the reversed arc of .
The adjacency matrix of a multidigraph with loops is the integer-valued matrix with rows and columns corresponding to the vertices, where a nondiagonal entry aij is the number of arcs from vertex i to vertex j, and the diagonal entry aii is the number of loops at vertex i. The adjacency matrix of a directed graph is a logical matrix, and is
unique up to permutation of rows and columns.
Another matrix representation for a directed graph is its incidence matrix.
See direction for more definitions.
Indegree and outdegree
For a vertex, the number of head ends adjacent to a vertex is called the indegree of the vertex and the number of tail ends adjacent to a vertex is its outdegree (called branching factor in trees).
Let and . The indegree of v is denoted deg−(v) and its outdegree is denoted deg+(v).
A vertex with is called a source, as it is the origin of each of its outgoing arcs. Similarly, a vertex with is called a sink, since it is the end of each of its incoming arcs.
The degree sum formula states that, for a directed graph,
If for every vertex , , the graph is called a balanced directed graph.
Degree sequence
The degree sequence of a directed graph is the list of its indegree and outdegree pairs; for the above example we have degree sequence ((2, 0), (2, 2), (0, 2), (1, 1)). The degree sequence is a directed graph invariant so isomorphic directed graphs have the same degree sequence. However, the degree sequence does not, in general, uniquely identify a directed graph; in some cases, non-isomorphic digraphs have the same degree sequence.
The directed graph realization problem is the problem of finding a directed graph with the degree sequence a given sequence of positive integer pairs. (Trailing pairs of zeros may be ignored since they are trivially realized by adding an appropriate number of isolated vertices to the directed graph.) A sequence which is the degree sequence of some directed graph, i.e. for which the directed graph realization problem has a solution, is called a directed graphic or directed graphical sequence. This problem can either be solved by the Kleitman–Wang algorithm or by the Fulkerson–Chen–Anstee theorem.
Directed graph connectivity
A directed graph is weakly connected (or just connected) if the undirected underlying graph obtained by replacing all directed edges of the graph with undirected edges is a connected graph.
A directed graph is strongly connected or strong if it contains a directed path from x to y (and from y to x) for every pair of vertices . The strong components are the maximal strongly connected subgraphs.
A connected rooted graph (or flow graph) is one where there exists a directed path to every vertex from a distinguished root vertex.
| Mathematics | Graph theory | null |
19725090 | https://en.wikipedia.org/wiki/Cold | Cold | Cold is the presence of low temperature, especially in the atmosphere. In common usage, cold is often a subjective perception. A lower bound to temperature is absolute zero, defined as 0.00K on the Kelvin scale, an absolute thermodynamic temperature scale. This corresponds to on the Celsius scale, on the Fahrenheit scale, and on the Rankine scale.
Since temperature relates to the thermal energy held by an object or a sample of matter, which is the kinetic energy of the random motion of the particle constituents of matter, an object will have less thermal energy when it is colder and more when it is hotter. If it were possible to cool a system to absolute zero, all motion of the particles in a sample of matter would cease and they would be at complete rest in the classical sense. The object could be described as having zero thermal energy. Microscopically in the description of quantum mechanics, however, matter still has zero-point energy even at absolute zero, because of the uncertainty principle.
Cooling
Cooling refers to the process of becoming cold, or lowering in temperature. This could be accomplished by removing heat from a system, or exposing the system to an environment with a lower temperature.
Coolants are fluids used to cool objects, prevent freezing and prevent erosion in machines.
Air cooling is the process of cooling an object by exposing it to air. This will only work if the air is at a lower temperature than the object, and the process can be enhanced by increasing the surface area, increasing the coolant flow rate, or decreasing the mass of the object.
Another common method of cooling is exposing an object to ice, dry ice, or liquid nitrogen. This works by conduction; the heat is transferred from the relatively warm object to the relatively cold coolant.
Laser cooling and magnetic evaporative cooling are techniques used to reach very low temperatures.
History
Early history
In ancient times, ice was not adopted for food preservation but used to cool wine which the Romans had also done. According to Pliny, Emperor Nero invented the ice bucket to chill wines instead of adding it to wine to make it cold as it would dilute it.
Some time around 1700 BC Zimri-Lim, king of Mari Kingdom in northwest Iraq had created an "icehouse" called bit shurpin at a location close to his capital city on the banks of the Euphrates. In the 7th century BC the Chinese had used icehouses to preserve vegetables and fruits. During the Tang dynastic rule in China (618–907 AD) a document refers to the practice of using ice that was in vogue during the Eastern Chou Dynasty (770–256 BC) by 94 workmen employed for "Ice-Service" to freeze everything from wine to dead bodies.
Shachtman says that in the 4th century AD, the brother of the Japanese emperor Nintoku gave him a gift of ice from a mountain. The Emperor was so happy with the gift that he named the first of June as the "Day of Ice" and ceremoniously gave blocks of ice to his officials.
Even in ancient times, Shachtman says, in Egypt and India, night cooling by evaporation of water and heat radiation, and the ability of salts to lower the freezing temperature of water was practiced. The ancient people of Rome and Greece were aware that boiled water cooled quicker than the ordinary water; the reason for this is that with boiling of water carbon dioxide and other gases, which are deterrents to cooling, are removed; but this fact was not known till the 17th century.
From the 17th century
Shachtman says that King James VI and I supported the work of Cornelis Drebbel as a magician to perform tricks such as producing thunder, lightning, lions, birds, trembling leaves and so forth. In 1620 he gave a demonstration in Westminster Abbey to the king and his courtiers on the power of cold. On a summer day, Shachtman says, Drebbel had created a chill (lowered the temperature by several degrees) in the hall of the Abbey, which made the king shiver and run out of the hall with his entourage. This was an incredible spectacle, says Shachtman. Several years before, Giambattista della Porta had demonstrated at the Abbey "ice fantasy gardens, intricate ice sculptures" and also iced drinks for banquets in Florence. The only reference to the artificial freezing created by Drebbel was by Francis Bacon. His demonstration was not taken seriously as it was considered one of his magic tricks, as there was no practical application then. Drebbel had not revealed his secrets.
Shachtman says that Lord Chancellor Bacon, an advocate of experimental science, had tried in Novum Organum, published in the late 1620s, to explain the artificial freezing experiment at Westminster Abbey, though he was not present during the demonstration, as "Nitre (or rather its spirit) is very cold, and hence nitre or salt when added to snow or ice intensifies the cold of the latter, the nitre by adding to its own cold, but the salt by supplying activity to the cold snow." This explanation on the cold inducing aspects of nitre and salt was tried then by many scientists.
Shachtman says it was the lack of scientific knowledge in physics and chemistry that had held back progress in the beneficial use of ice until a drastic change in religious opinions in the 17th century. The intellectual barrier was broken by Francis Bacon and Robert Boyle who followed him in this quest for knowledge of cold. Boyle did extensive experimentation during the 17th century in the discipline of cold, and his research on pressure and volume was the forerunner of research in the field of cold during the 19th century. He explained his approach as "Bacon's identification of heat and cold as the right and left hands of nature". Boyle also refuted some of the theories mooted by Aristotle on cold by experimenting on transmission of cold from one material to the other. He proved that water was not the only source of cold but gold, silver and crystal, which had no water content, could also change to severe cold condition.
19th century
In the United States from about 1850 till end of 19th century export of ice was second only to cotton. The first ice box was developed by Thomas Moore, a farmer from Maryland in 1810 to carry butter in an oval shaped wooden tub. The tub was provided with a metal lining in its interior and surrounded by a packing of ice. A rabbit skin was used as insulation. Moore also developed an ice box for domestic use with the container built over a space of which was filled with ice. In 1825, ice harvesting by use of a horse drawn ice cutting device was invented by Nathaniel J. Wyeth. The cut blocks of uniform size ice was a cheap method of food preservation widely practiced in the United States. Also developed in 1855 was a steam powered device to haul 600 tons of ice per hour. More innovations ensued. Devices using compressed air as a refrigerants were invented.
20th century
Iceboxes were in widespread use from the mid-19th century to the 1930s, when the refrigerator was introduced into the home. Most municipally consumed ice was harvested in winter from snow-packed areas or frozen lakes, stored in ice houses, and delivered domestically as iceboxes became more common.
In 1913, refrigerators for home use were invented. In 1923 Frigidaire introduced the first self-contained unit. The introduction of Freon in the 1920s expanded the refrigerator market during the 1930s. Home freezers as separate compartments (larger than necessary just for ice cubes) were introduced in 1940. Frozen foods, previously a luxury item, became commonplace.
Physiological effects
Cold has numerous physiological and pathological effects on the human body, as well as on other organisms. Cold environments may promote certain psychological traits, as well as having direct effects on the ability to move. Shivering is one of the first physiological responses to cold. Even at low temperatures, the cold can massively disrupt blood circulation. Extracellular water freezes and tissue is destroyed. It affects fingers, toes, nose, ears and cheeks particularly often. They discolor, swell, blister, and bleed. The so-called frostnip leads to local frostbite or even to the death of entire body parts. Only temporary cold reactions of the skin are without consequences. As blood vessels contract, they become cool and pale, with less oxygen getting into the tissue. Warmth stimulates blood circulation again and is painful but harmless. Comprehensive protection against the cold is particularly important for children and for sports. Extreme cold temperatures may lead to frostbite, sepsis, and hypothermia, which in turn may result in death.
Common myths
A common, but false, statement states that cold weather itself can induce the identically named common cold. No scientific evidence of this has been found, although the disease, alongside influenza and others, does increase in prevalence with colder weather.
Notable cold locations and objects
The National Institute of Standards and Technology in Boulder, Colorado using a new technique, managed to chill a microscopic mechanical drum to 360 microkelvins, making it the coldest object on record. Theoretically, using this technique, an object could be cooled to absolute zero.
The coldest known temperature ever achieved is a state of matter called the Bose–Einstein condensate which was first theorized to exist by Satyendra Nath Bose in 1924 and first created by Eric Cornell, Carl Wieman, and co-workers at JILA on 5 June 1995. They did this by cooling a dilute vapor consisting of approximately two thousand rubidium-87 atoms to below 170 nK (one nK or nanokelvin is a billionth (10−9) of a kelvin) using a combination of laser cooling (a technique that won its inventors Steven Chu, Claude Cohen-Tannoudji, and William D. Phillips the 1997 Nobel Prize in Physics) and magnetic evaporative cooling.
90377 Sedna is one of the coldest known objects within the Solar System. Orbiting at an average distance of 84 billion miles, Sedna has an average surface temperature of -400°F (-240°C).
The lunar crater Hermite was described after a 2009 survey by NASA's Lunar Reconnaissance Orbiter as the "coldest known place in the Solar System", with temperatures at 26 kelvins (−413 °F, −247 °C).
The Boomerang Nebula is the coldest known natural location in the universe, with a temperature that is estimated at 1 K (−272.15 °C, −457.87 °F).
The Dwarf Planet Haumea is one of the coldest known objects in our solar system. With a Temperature of -401 degrees Fahrenheit or -241 degrees Celsius
The Planck spacecraft's instruments are kept at 0.1 K (−273.05 °C, −459.49 °F) via passive and active cooling.
Absent any other source of heat, the temperature of the Universe is roughly 2.725 kelvins, due to the Cosmic microwave background radiation, a remnant of the Big Bang.
Neptune's moon Triton has a surface temperature of 38.15 K (−235 °C, −391 °F)
Uranus with a black-body temperature of 58.2 K (−215.0 °C, −354.9 °F).
Saturn with a black-body temperature of 81.1 K (−192.0 °C, −313.7 °F).
Mercury, despite being close to the Sun, is actually cold during its night, with a temperature of about 93.15 K (−180 °C, −290 °F). Mercury is cold during its night because it has no atmosphere to trap in heat from the Sun.
Jupiter with a black-body temperature of 110.0 K (−163.2 °C, −261.67 °F).
Mars with a black-body temperature of 210.1 K (−63.05 °C, −81.49 °F).
The coldest continent on Earth is Antarctica. The coldest place on Earth is the Antarctic Plateau, an area of Antarctica around the South Pole that has an altitude of around . The lowest reliably measured temperature on Earth of 183.9 K (−89.2 °C, −128.6 °F) was recorded there at Vostok Station on 21 July 1983. The Poles of Cold are the places in the Southern and Northern Hemispheres where the lowest air temperatures have been recorded. (See List of weather records).
The cold deserts of the North Pole, known as the tundra region, experiences an annual snow fall of a few inches and temperatures recorded are as low as 203.15 K (−70 °C, −94 °F). Only a few small plants survive in the generally frozen ground (thaws only for a short spell).
Cold deserts of the Himalayas are a feature of a rain-shadow zone created by the mountain peaks of the Himalaya range that runs from Pamir Knot extending to the southern border of the Tibetan plateau; however this mountain range is also the reason for the monsoon rain fall in the Indian subcontinent. This zone is located in an elevation of about 3,000 m, and covers Ladakh, Lahaul, Spiti and Pooh. In addition, there are inner valleys within the main Himalayas such as Chamoli, some areas of Kinnaur, Pithoragarh and northern Sikkim which are also categorized as cold deserts.
Mythology and culture
Niflheim was a realm of primordial ice and cold with nine frozen rivers in Norse Mythology.
The "Hell in Dante's Inferno" is stated as Cocytus a frozen lake where Virgil and Dante were deposited.
| Physical sciences | Thermodynamics | Physics |
19727024 | https://en.wikipedia.org/wiki/Rational%20number | Rational number | In mathematics, a rational number is a number that can be expressed as the quotient or fraction of two integers, a numerator and a non-zero denominator . For example, is a rational number, as is every integer (for example, The set of all rational numbers, also referred to as "the rationals", the field of rationals or the field of rational numbers is usually denoted by boldface , or blackboard bold
A rational number is a real number. The real numbers that are rational are those whose decimal expansion either terminates after a finite number of digits (example: ), or eventually begins to repeat the same finite sequence of digits over and over (example: ). This statement is true not only in base 10, but also in every other integer base, such as the binary and hexadecimal ones (see ).
A real number that is not rational is called irrational. Irrational numbers include the square root of 2 , , and the golden ratio (). Since the set of rational numbers is countable, and the set of real numbers is uncountable, almost all real numbers are irrational.
Rational numbers can be formally defined as equivalence classes of pairs of integers with , using the equivalence relation defined as follows:
The fraction then denotes the equivalence class of .
Rational numbers together with addition and multiplication form a field which contains the integers, and is contained in any field containing the integers. In other words, the field of rational numbers is a prime field, and a field has characteristic zero if and only if it contains the rational numbers as a subfield. Finite extensions of are called algebraic number fields, and the algebraic closure of is the field of algebraic numbers.
In mathematical analysis, the rational numbers form a dense subset of the real numbers. The real numbers can be constructed from the rational numbers by completion, using Cauchy sequences, Dedekind cuts, or infinite decimals (see Construction of the real numbers).
Terminology
In mathematics, "rational" is often used as a noun abbreviating "rational number". The adjective rational sometimes means that the coefficients are rational numbers. For example, a rational point is a point with rational coordinates (i.e., a point whose coordinates are rational numbers); a rational matrix is a matrix of rational numbers; a rational polynomial may be a polynomial with rational coefficients, although the term "polynomial over the rationals" is generally preferred, to avoid confusion between "rational expression" and "rational function" (a polynomial is a rational expression and defines a rational function, even if its coefficients are not rational numbers). However, a rational curve is not a curve defined over the rationals, but a curve which can be parameterized by rational functions.
Etymology
Although nowadays rational numbers are defined in terms of ratios, the term rational is not a derivation of ratio. On the contrary, it is ratio that is derived from rational: the first use of ratio with its modern meaning was attested in English about 1660, while the use of rational for qualifying numbers appeared almost a century earlier, in 1570. This meaning of rational came from the mathematical meaning of irrational, which was first used in 1551, and it was used in "translations of Euclid (following his peculiar use of )".
This unusual history originated in the fact that ancient Greeks "avoided heresy by forbidding themselves from thinking of those [irrational] lengths as numbers". So such lengths were irrational, in the sense of illogical, that is "not to be spoken about" ( in Greek).
Arithmetic
Irreducible fraction
Every rational number may be expressed in a unique way as an irreducible fraction where and are coprime integers and . This is often called the canonical form of the rational number.
Starting from a rational number its canonical form may be obtained by dividing and by their greatest common divisor, and, if , changing the sign of the resulting numerator and denominator.
Embedding of integers
Any integer can be expressed as the rational number which is its canonical form as a rational number.
Equality
if and only if
If both fractions are in canonical form, then:
if and only if and
Ordering
If both denominators are positive (particularly if both fractions are in canonical form):
if and only if
On the other hand, if either denominator is negative, then each fraction with a negative denominator must first be converted into an equivalent form with a positive denominator—by changing the signs of both its numerator and denominator.
Addition
Two fractions are added as follows:
If both fractions are in canonical form, the result is in canonical form if and only if are coprime integers.
Subtraction
If both fractions are in canonical form, the result is in canonical form if and only if are coprime integers.
Multiplication
The rule for multiplication is:
where the result may be a reducible fraction—even if both original fractions are in canonical form.
Inverse
Every rational number has an additive inverse, often called its opposite,
If is in canonical form, the same is true for its opposite.
A nonzero rational number has a multiplicative inverse, also called its reciprocal,
If is in canonical form, then the canonical form of its reciprocal is either or depending on the sign of .
Division
If are nonzero, the division rule is
Thus, dividing by is equivalent to multiplying by the reciprocal of
Exponentiation to integer power
If is a non-negative integer, then
The result is in canonical form if the same is true for In particular,
If , then
If is in canonical form, the canonical form of the result is if or is even. Otherwise, the canonical form of the result is
Continued fraction representation
A finite continued fraction is an expression such as
where are integers. Every rational number can be represented as a finite continued fraction, whose coefficients can be determined by applying the Euclidean algorithm to .
Other representations
common fraction:
mixed numeral:
repeating decimal using a vinculum:
repeating decimal using parentheses:
continued fraction using traditional typography:
continued fraction in abbreviated notation:
Egyptian fraction:
prime power decomposition:
quote notation:
are different ways to represent the same rational value.
Formal construction
The rational numbers may be built as equivalence classes of ordered pairs of integers.
More precisely, let be the set of the pairs of integers such . An equivalence relation is defined on this set by
Addition and multiplication can be defined by the following rules:
This equivalence relation is a congruence relation, which means that it is compatible with the addition and multiplication defined above; the set of rational numbers is the defined as the quotient set by this equivalence relation, equipped with the addition and the multiplication induced by the above operations. (This construction can be carried out with any integral domain and produces its field of fractions.)
The equivalence class of a pair is denoted
Two pairs and belong to the same equivalence class (that is are equivalent) if and only if
This means that
if and only if
Every equivalence class may be represented by infinitely many pairs, since
Each equivalence class contains a unique canonical representative element. The canonical representative is the unique pair in the equivalence class such that and are coprime, and . It is called the representation in lowest terms of the rational number.
The integers may be considered to be rational numbers identifying the integer with the rational number
A total order may be defined on the rational numbers, that extends the natural order of the integers. One has
If
Properties
The set of all rational numbers, together with the addition and multiplication operations shown above, forms a field.
has no field automorphism other than the identity. (A field automorphism must fix 0 and 1; as it must fix the sum and the difference of two fixed elements, it must fix every integer; as it must fix the quotient of two fixed elements, it must fix every rational number, and is thus the identity.)
is a prime field, which is a field that has no subfield other than itself. The rationals are the smallest field with characteristic zero. Every field of characteristic zero contains a unique subfield isomorphic to
With the order defined above, is an ordered field that has no subfield other than itself, and is the smallest ordered field, in the sense that every ordered field contains a unique subfield isomorphic to
is the field of fractions of the integers The algebraic closure of i.e. the field of roots of rational polynomials, is the field of algebraic numbers.
The rationals are a densely ordered set: between any two rationals, there sits another one, and, therefore, infinitely many other ones. For example, for any two fractions such that
(where are positive), we have
Any totally ordered set which is countable, dense (in the above sense), and has no least or greatest element is order isomorphic to the rational numbers.
Countability
The set of all rational numbers is countable, as is illustrated in the figure to the right. As a rational number can be expressed as a ratio of two integers, it is possible to assign two integers to any point on a square lattice as in a Cartesian coordinate system, such that any grid point corresponds to a rational number. This method, however, exhibits a form of redundancy, as several different grid points will correspond to the same rational number; these are highlighted in red on the provided graphic. An obvious example can be seen in the line going diagonally towards the bottom right; such ratios will always equal 1, as any non-zero number divided by itself will always equal one.
It is possible to generate all of the rational numbers without such redundancies: examples include the Calkin–Wilf tree and Stern–Brocot tree.
As the set of all rational numbers is countable, and the set of all real numbers (as well as the set of irrational numbers) is uncountable, the set of rational numbers is a null set, that is, almost all real numbers are irrational, in the sense of Lebesgue measure.
Real numbers and topological properties
The rationals are a dense subset of the real numbers; every real number has rational numbers arbitrarily close to it. A related property is that rational numbers are the only numbers with finite expansions as regular continued fractions.
In the usual topology of the real numbers, the rationals are neither an open set nor a closed set.
By virtue of their order, the rationals carry an order topology. The rational numbers, as a subspace of the real numbers, also carry a subspace topology. The rational numbers form a metric space by using the absolute difference metric and this yields a third topology on All three topologies coincide and turn the rationals into a topological field. The rational numbers are an important example of a space which is not locally compact. The rationals are characterized topologically as the unique countable metrizable space without isolated points. The space is also totally disconnected. The rational numbers do not form a complete metric space, and the real numbers are the completion of under the metric above.
p-adic numbers
In addition to the absolute value metric mentioned above, there are other metrics which turn into a topological field:
Let be a prime number and for any non-zero integer , let where is the highest power of dividing .
In addition set For any rational number we set
Then
defines a metric on
The metric space is not complete, and its completion is the -adic number field Ostrowski's theorem states that any non-trivial absolute value on the rational numbers is equivalent to either the usual real absolute value or a -adic absolute value.
| Mathematics | Basics | null |
3559747 | https://en.wikipedia.org/wiki/Staggered%20conformation | Staggered conformation | In organic chemistry, a staggered conformation is a chemical conformation of an ethane-like moiety abcX–Ydef in which the substituents a, b, and c are at the maximum distance from d, e, and f; this requires the torsion angles to be 60°. It is the opposite of an eclipsed conformation, in which those substituents are as close to each other as possible.
Such a conformation exists in any open chain single chemical bond connecting two sp3-hybridised atoms, and is normally a conformational energy minimum. For some molecules such as those of n-butane, there can be special versions of staggered conformations called gauche and anti; see first Newman projection diagram in conformational isomerism.
Staggered/eclipsed configurations also distinguish different crystalline structures of e.g. cubic/hexagonal boron nitride, and diamond/lonsdaleite.
| Physical sciences | Stereochemistry | Chemistry |
3561094 | https://en.wikipedia.org/wiki/Modifiable%20areal%20unit%20problem | Modifiable areal unit problem |
The modifiable areal unit problem (MAUP) is a source of statistical bias that can significantly impact the results of statistical hypothesis tests. MAUP affects results when point-based measures of spatial phenomena are aggregated into spatial partitions or areal units (such as regions or districts) as in, for example, population density or illness rates. The resulting summary values (e.g., totals, rates, proportions, densities) are influenced by both the shape and scale of the aggregation unit.
For example, census data may be aggregated into county districts, census tracts, postcode areas, police precincts, or any other arbitrary spatial partition. Thus the results of data aggregation are dependent on the mapmaker's choice of which "modifiable areal unit" to use in their analysis. A census choropleth map calculating population density using state boundaries will yield radically different results than a map that calculates density based on county boundaries. Furthermore, census district boundaries are also subject to change over time, meaning the MAUP must be considered when comparing past data to current data.
Background
The issue was first recognized by Gehlke and Biehl in 1934 and later described in detail in an entry in the Concepts and Techniques in Modern Geography (CATMOG) series by Stan Openshaw (1984) and in the book by Giuseppe Arbia (1988). In particular, Openshaw (1984) observed that "the areal units (zonal objects) used in many geographical studies are arbitrary, modifiable, and subject to the whims and fancies of whoever is doing, or did, the aggregating". The problem is especially apparent when the aggregate data are used for cluster analysis for spatial epidemiology, spatial statistics or choropleth mapping, in which misinterpretations can easily be made without realizing it. Many fields of science, especially human geography are prone to disregard the MAUP when drawing inferences from statistics based on aggregated data. MAUP is closely related to the topic of ecological fallacy and ecological bias (Arbia, 1988). Stan Openshaw's work on this topic has led to Michael F. Goodchild suggesting it be referred to as the "Openshaw effect."
Ecological bias caused by MAUP has been documented as two separate effects that usually occur simultaneously during the analysis of aggregated data. First, the scale effect causes variation in statistical results between different levels of aggregation (radial distance). Therefore, the association between variables depends on the size of areal units for which data are reported. Generally, correlation increases as areal unit size increases. The zoning effect describes variation in correlation statistics caused by the regrouping of data into different configurations at the same scale (areal shape).
Since the 1930s, research has found extra variation in statistical results because of the MAUP. The standard methods of calculating within-group and between-group variance do not account for the extra variance seen in MAUP studies as the groupings change. MAUP can be used as a methodology to calculate upper and lower limits as well as average regression parameters for multiple sets of spatial groupings. The MAUP is a critical source of error in spatial studies, whether observational or experimental. As such, unit consistency, particularly in a time-series cross-sectional (TSCS) context, is essential. Further, robustness checks of unit sensitivity to alternative spatial aggregation should be routinely performed to mitigate associated biases on resulting statistical estimates.
Suggested solutions
Several suggestions have been made in literature to reduce aggregation bias during regression analysis. A researcher might correct the variance-covariance matrix using samples from individual-level data. Alternatively, one might focus on local spatial regression rather than global regression. A researcher might also attempt to design areal units to maximize a particular statistical result. Others have argued that it may be difficult to construct a single set of optimal aggregation units for multiple variables, each of which may exhibit non-stationarity and spatial autocorrelation across space in different ways. Others have suggested developing statistics that change across scales in a predictable way, perhaps using fractal dimension as a scale-independent measure of spatial relationships. Others have suggested Bayesian hierarchical models as a general methodology for combining aggregated and individual-level data for ecological inference.
Studies of the MAUP based on empirical data can only provide limited insight due to an inability to control relationships between multiple spatial variables. Data simulation is necessary to have control over various properties of individual-level data. Simulation studies have demonstrated that the spatial support of variables can affect the magnitude of ecological bias caused by spatial data aggregation.
MAUP sensitivity analysis
Using simulations for univariate data, Larsen advocated the use of a Variance Ratio to investigate the effect of spatial configuration, spatial association, and data aggregation. A detailed description of the variation of statistics due to MAUP is presented by Reynolds, who demonstrates the importance of the spatial arrangement and spatial autocorrelation of data values. Reynold’s simulation experiments were expanded by Swift, who in which a series of nine exercises began with simulated regression analysis and spatial trend, then focused on the topic of MAUP in the context of spatial epidemiology. A method of MAUP sensitivity analysis is presented that demonstrates that the MAUP is not entirely a problem. MAUP can be used as an analytical tool to help understand spatial heterogeneity and spatial autocorrelation.
This topic is of particular importance because in some cases data aggregation can obscure a strong correlation between variables, making the relationship appear weak or even negative. Conversely, MAUP can cause random variables to appear as if there is a significant association where there is not. Multivariate regression parameters are more sensitive to MAUP than correlation coefficients. Until a more analytical solution to MAUP is discovered, spatial sensitivity analysis using a variety of areal units is recommended as a methodology to estimate the uncertainty of correlation and regression coefficients due to ecological bias. An example of data simulation and re-aggregation using the ArcPy library is available.
In transport planning, MAUP is associated to Traffic Analysis Zoning (TAZ). A major point of departure in understanding problems in transportation analysis is the recognition that spatial analysis has some limitations associated with the discretization of space. Among them, modifiable areal units and boundary problems are directly or indirectly related to transportation planning and analysis through the design of traffic analysis zones – most of transport studies require directly or indirectly the definition of TAZs. The modifiable boundary and the scale issues should all be given specific attention during the specification of a TAZ because of the effects these factors exert on statistical and mathematical properties of spatial patterns (ie the modifiable areal unit problem—MAUP). In the studies of Viegas, Martinez and Silva (2009, 2009b) the authors propose a method where the results obtained from the study of spatial data are not independent of the scale, and the aggregation effects are implicit in the choice of zonal boundaries. The delineation of zonal boundaries of TAZs has a direct impact on the reality and accuracy of the results obtained from transportation forecasting models. In this paper the MAUP effects on the TAZ definition and the transportation demand models are measured and analyzed using different grids (in size and in origin location). This analysis was developed by building an application integrated in commercial GIS software and by using a case study (Lisbon Metropolitan Area) to test its implementabiity and performance. The results reveal the conflict between statistical and geographic precision, and their relationship with the loss of information in the traffic assignment step of the transportation planning models.
Research has also identified the modifiable areal unit problem (MAUP) to be a factor in climate action and governance by affecting coordination between national and local actors. Data scaling issues associated with MAUP may result in mismatches in climate priorities and create inequities in the outcomes of climate action, potentially undermining the effectiveness of policies designed to address climate change at different governance levels.
| Mathematics | Statistics | null |
125769 | https://en.wikipedia.org/wiki/Binding%20energy | Binding energy | In physics and chemistry, binding energy is the smallest amount of energy required to remove a particle from a system of particles or to disassemble a system of particles into individual parts. In the former meaning the term is predominantly used in condensed matter physics, atomic physics, and chemistry, whereas in nuclear physics the term separation energy is used. A bound system is typically at a lower energy level than its unbound constituents. According to relativity theory, a decrease in the total energy of a system is accompanied by a decrease in the total mass, where .
Types
There are several types of binding energy, each operating over a different distance and energy scale. The smaller the size of a bound system, the higher its associated binding energy.
Mass–energy relation
A bound system is typically at a lower energy level than its unbound constituents because its mass must be less than the total mass of its unbound constituents. For systems with low binding energies, this "lost" mass after binding may be fractionally small, whereas for systems with high binding energies, the missing mass may be an easily measurable fraction. This missing mass may be lost during the process of binding as energy in the form of heat or light, with the removed energy corresponding to the removed mass through Einstein's equation . In the process of binding, the constituents of the system might enter higher energy states of the nucleus/atom/molecule while retaining their mass, and because of this, it is necessary that they are removed from the system before its mass can decrease. Once the system cools to normal temperatures and returns to ground states regarding energy levels, it will contain less mass than when it first combined and was at high energy. This loss of heat represents the "mass deficit", and the heat itself retains the mass that was lost (from the point of view of the initial system). This mass will appear in any other system that absorbs the heat and gains thermal energy.
For example, if two objects are attracting each other in space through their gravitational field, the attraction force accelerates the objects, increasing their velocity, which converts their potential energy (gravity) into kinetic energy. When the particles either pass through each other without interaction or elastically repel during the collision, the gained kinetic energy (related to speed) begins to revert into potential energy, driving the collided particles apart. The decelerating particles will return to the initial distance and beyond into infinity, or stop and repeat the collision (oscillation takes place). This shows that the system, which loses no energy, does not combine (bind) into a solid object, parts of which oscillate at short distances. Therefore, to bind the particles, the kinetic energy gained due to the attraction must be dissipated by resistive force. Complex objects in collision ordinarily undergo inelastic collision, transforming some kinetic energy into internal energy (heat content, which is atomic movement), which is further radiated in the form of photonsthe light and heat. Once the energy to escape the gravity is dissipated in the collision, the parts will oscillate at a closer, possibly atomic, distance, thus looking like one solid object. This lost energy, necessary to overcome the potential barrier to separate the objects, is the binding energy. If this binding energy were retained in the system as heat, its mass would not decrease, whereas binding energy lost from the system as heat radiation would itself have mass. It directly represents the "mass deficit" of the cold, bound system.
Closely analogous considerations apply in chemical and nuclear reactions. Exothermic chemical reactions in closed systems do not change mass, but do become less massive once the heat of reaction is removed, though this mass change is too small to measure with standard equipment. In nuclear reactions, the fraction of mass that may be removed as light or heat, i.e. binding energy, is often a much larger fraction of the system mass. It may thus be measured directly as a mass difference between rest masses of reactants and (cooled) products. This is because nuclear forces are comparatively stronger than the Coulombic forces associated with the interactions between electrons and protons that generate heat in chemistry.
Mass change
Mass change (decrease) in bound systems, particularly atomic nuclei, has also been termed mass defect, mass deficit, or mass packing fraction.
The difference between the unbound system calculated mass and experimentally measured mass of nucleus (mass change) is denoted as Δm. It can be calculated as follows:
Mass change = (unbound system calculated mass) − (measured mass of system)
e.g. (sum of masses of protons and neutrons) − (measured mass of nucleus)
After a nuclear reaction occurs that results in an excited nucleus, the energy that must be radiated or otherwise removed as binding energy in order to decay to the unexcited state may be in one of several forms. This may be electromagnetic waves, such as gamma radiation; the kinetic energy of an ejected particle, such as an electron, in internal conversion decay; or partly as the rest mass of one or more emitted particles, such as the particles of beta decay. No mass deficit can appear, in theory, until this radiation or this energy has been emitted and is no longer part of the system.
When nucleons bind together to form a nucleus, they must lose a small amount of mass, i.e. there is a change in mass to stay bound. This mass change must be released as various types of photon or other particle energy as above, according to the relation . Thus, after the binding energy has been removed, binding energy = mass change × . This energy is a measure of the forces that hold the nucleons together. It represents energy that must be resupplied from the environment for the nucleus to be broken up into individual nucleons.
For example, an atom of deuterium has a mass defect of 0.0023884 Da, and its binding energy is nearly equal to 2.23 MeV. This means that energy of 2.23 MeV is required to disintegrate an atom of deuterium.
The energy given off during either nuclear fusion or nuclear fission is the difference of the binding energies of the "fuel", i.e. the initial nuclide(s), from that of the fission or fusion products. In practice, this energy may also be calculated from the substantial mass differences between the fuel and products, which uses previous measurements of the atomic masses of known nuclides, which always have the same mass for each species. This mass difference appears once evolved heat and radiation have been removed, which is required for measuring the (rest) masses of the (non-excited) nuclides involved in such calculations.
| Physical sciences | Atomic physics | null |
2614714 | https://en.wikipedia.org/wiki/Vertigo | Vertigo | Vertigo is a condition in which a person has the sensation that they are moving, or that objects around them are moving, when they are not. Often it feels like a spinning or swaying movement. It may be associated with nausea, vomiting, perspiration, or difficulties walking. It is typically worse when the head is moved. Vertigo is the most common type of dizziness.
The most common disorders that result in vertigo are benign paroxysmal positional vertigo (BPPV), Ménière's disease, and vestibular neuritis. Less common causes include stroke, brain tumors, brain injury, multiple sclerosis, migraines, trauma, and uneven pressures between the middle ears. Physiologic vertigo may occur following being exposed to motion for a prolonged period such as when on a ship or simply following spinning with the eyes closed. Other causes may include toxin exposures such as to carbon monoxide, alcohol, or aspirin. Vertigo typically indicates a problem in a part of the vestibular system. Other causes of dizziness include presyncope, disequilibrium, and non-specific dizziness.
Benign paroxysmal positional vertigo is more likely in someone who gets repeated episodes of vertigo with movement and is otherwise normal between these episodes. Benign vertigo episodes generally last less than one minute. The Dix-Hallpike test typically produces a period of rapid eye movements known as nystagmus in this condition. In Ménière's disease there is often ringing in the ears, hearing loss, and the attacks of vertigo last more than twenty minutes. In vestibular neuritis the onset of vertigo is sudden, and the nystagmus occurs even when the person has not been moving. In this condition vertigo can last for days. More severe causes should also be considered, especially if other problems such as weakness, headache, double vision, or numbness occur.
Dizziness affects approximately 20–40% of people at some point in time, while about 7.5–10% have vertigo. About 5% have vertigo in a given year. It becomes more common with age and affects women two to three times more often than men. Vertigo accounts for about 2–3% of emergency department visits in the developed world.
Classification
Vertigo is classified into either peripheral or central depending on the location of the dysfunction of the vestibular pathway, although it can also be caused by psychological factors.
Vertigo can also be classified into objective, subjective, and pseudovertigo. Objective vertigo describes when the person has the sensation that stationary objects in the environment are moving. Subjective vertigo refers to when the person feels as if they are moving. The third type is known as pseudovertigo, an intensive sensation of rotation inside the person's head. While this classification appears in textbooks, it is unclear what relation it has to the pathophysiology or treatment of vertigo.
Peripheral
Vertigo that is caused by problems with the inner ear or vestibular system, which is composed of the semicircular canals, the vestibule (utricle and saccule), and the vestibular nerve is called "peripheral", "otologic", or "vestibular" vertigo. The most common cause is benign paroxysmal positional vertigo (BPPV), which accounts for 32% of all peripheral vertigo. Other causes include Ménière's disease (12%), superior canal dehiscence syndrome, vestibular neuritis, and visual vertigo. Any cause of inflammation such as common cold, influenza, and bacterial infections may cause transient vertigo if it involves the inner ear, as may chemical insults (e.g., aminoglycosides) or physical trauma (e.g., skull fractures). Motion sickness is sometimes classified as a cause of peripheral vertigo.
People with peripheral vertigo typically present with mild to moderate imbalance, nausea, vomiting, hearing loss, tinnitus, fullness, and pain in the ear. In addition, lesions of the internal auditory canal may be associated with facial weakness on the same side. Due to a rapid compensation process, acute vertigo as a result of a peripheral lesion tends to improve in a short period of time (days to weeks).
Central
Vertigo that arises from injury to the balance centers of the central nervous system (CNS), often from a lesion in the brainstem or cerebellum, is called "central" vertigo and is generally associated with less prominent movement illusion and nausea than vertigo of peripheral origin. Central vertigo may have accompanying neurologic deficits (such as slurred speech and double vision), and pathologic nystagmus (which is pure vertical/torsional). Central pathology can cause disequilibrium, which is the sensation of being off balance. The balance disorder associated with central lesions causing vertigo is often so severe that many people are unable to stand or walk.
A number of conditions that involve the central nervous system may lead to vertigo including: lesions caused by infarctions or hemorrhage, tumors present in the cerebellopontine angle such as a vestibular schwannoma or cerebellar tumors, epilepsy, cervical spine disorders such as cervical spondylosis, degenerative ataxia disorders, migraine headaches, lateral medullary syndrome, Chiari malformation, multiple sclerosis, parkinsonism, as well as cerebral dysfunction. Central vertigo may not improve or may do so more slowly than vertigo caused by disturbance to peripheral structures. Alcohol can result in positional alcohol nystagmus (PAN).
Signs and symptoms
Vertigo is a sensation of spinning while stationary. It is commonly associated with nausea or vomiting, unsteadiness (postural instability), falls, changes to a person's thoughts, and difficulties in walking. Recurrent episodes in those with vertigo are common and frequently impair the quality of life. Blurred vision, difficulty in speaking, a lowered level of consciousness, and hearing loss may also occur. The signs and symptoms of vertigo can present as a persistent (insidious) onset or an episodic (sudden) onset.
Persistent onset vertigo is characterized by symptoms lasting for longer than one day and is caused by degenerative changes that affect balance as people age. Nerve conduction slows with aging, and a decreased vibratory sensation is common as a result. Additionally, there is a degeneration of the ampulla and otolith organs with an increase in age. Persistent onset is commonly paired with central vertigo signs and symptoms.
The characteristics of an episodic onset vertigo are indicated by symptoms lasting for a smaller, more memorable amount of time, typically lasting for only seconds to minutes.
Pathophysiology
The neurochemistry of vertigo includes six primary neurotransmitters that have been identified between the three-neuron arc that drives the vestibulo-ocular reflex (VOR). Glutamate maintains the resting discharge of the central vestibular neurons and may modulate synaptic transmission in all three neurons of the VOR arc. Acetylcholine appears to function as an excitatory neurotransmitter in both the peripheral and central synapses. Gamma-Aminobutyric acid (GABA) is thought to be inhibitory for the commissures of the medial vestibular nucleus, the connections among the cerebellar Purkinje cells, the lateral vestibular nucleus, and the vertical VOR.
Three other neurotransmitters work centrally. Dopamine may accelerate vestibular compensation. Norepinephrine modulates the intensity of central reactions to vestibular stimulation and facilitates compensation. Histamine is present only centrally, but its role is unclear. Dopamine, histamine, serotonin, and acetylcholine are neurotransmitters thought to produce vomiting. It is known that centrally acting antihistamines modulate the symptoms of acute symptomatic vertigo.
Diagnosis
Tests for vertigo often attempt to elicit nystagmus and to differentiate vertigo from other causes of dizziness such as presyncope, hyperventilation syndrome, disequilibrium, or psychiatric causes of lightheadedness. Tests of vestibular system (balance) function include electronystagmography (ENG), Dix-Hallpike maneuver, rotation tests, head-thrust test, caloric reflex test, and computerized dynamic posturography (CDP).
The HINTS test, which is a combination of three physical examination tests that may be performed by physicians at the bedside, has been deemed helpful in differentiating between central and peripheral causes of vertigo. The HINTS test involves the horizontal head impulse test, observation of nystagmus on primary gaze, and the test of skew. CT scans or MRIs are sometimes used by physicians when diagnosing vertigo.
Tests of auditory system (hearing) function include pure tone audiometry, speech audiometry, acoustic reflex, electrocochleography (ECoG), otoacoustic emissions (OAE), and the auditory brainstem response test.
A number of specific conditions can cause vertigo. In the elderly, however, the condition is often multifactorial.
A recent history of underwater diving can indicate a possibility of barotrauma or decompression sickness involvement, but does not exclude all other possibilities. The dive profile (which is frequently recorded by dive computer) can be useful to assess a probability for decompression sickness, which can be confirmed by therapeutic recompression.
Benign paroxysmal positional vertigo
Benign paroxysmal positional vertigo (BPPV) is the most common vestibular disorder and occurs when loose calcium carbonate debris has broken off of the otoconial membrane and enters a semicircular canal thereby creating the sensation of motion. People with BPPV may experience brief periods of vertigo, usually under a minute, which occur with change in the position.
This is the most common cause of vertigo. It occurs in 0.6% of the population yearly with 10% having an attack during their lifetime. It is believed to be due to a mechanical malfunction of the inner ear. BPPV may be diagnosed with the Dix-Hallpike test and can be effectively treated with repositioning movements such as the Epley maneuver.
Ménière's disease
Ménière's disease is an inner ear disorder of unknown origin, but is thought to be caused by an increase in the amount of endolymphatic fluid present in the inner ear (endolymphatic hydrops). However, this idea has not been directly confirmed with histopathologic studies, but electrophysiologic studies have been suggestive of this mechanism. Ménière's disease frequently presents with recurrent, spontaneous attacks of severe vertigo in combination with ringing in the ears (tinnitus), a feeling of pressure or fullness in the ear (aural fullness), severe nausea or vomiting, imbalance, and hearing loss. As the disease worsens, hearing loss will progress.
Vestibular neuritis
Vestibular neuritis presents with severe vertigo with associated nausea, vomiting, and generalized imbalance and is believed to be caused by a viral infection of the inner ear, although several theories have been put forward and the cause remains uncertain. Individuals with vestibular neuritis do not typically have auditory symptoms, but may experience a sensation of aural fullness or tinnitus. Persisting balance problems may remain in 30% of people affected.
Vestibular migraine
Vestibular migraine is the association of vertigo and migraines and is one of the most common causes of recurrent, spontaneous episodes of vertigo. The cause of vestibular migraines is currently unclear; however, one hypothesized cause is that the stimulation of the trigeminal nerve leads to nystagmus in individuals with migraines. Approximately 40% of all migraine patients will have an accompanying vestibular syndrome, such as vertigo, dizziness, or disruption of the balance system.
Other suggested causes of vestibular migraines include the following: unilateral neuronal instability of the vestibular nerve, idiopathic asymmetric activation of the vestibular nuclei in the brainstem, and vasospasm of the blood vessels supplying the labyrinth or central vestibular pathways resulting in ischemia to these structures. Vestibular migraines are estimated to affect 1–3% of the general population and may affect 10% of people with migraine . Additionally, vestibular migraines tend to occur more often in women and rarely affect individuals after the sixth decade of life.
Motion sickness
Motion sickness is common and is related to vestibular migraine. It is nausea and vomiting in response to motion and is typically worse if the journey is on a winding road or involves many stops and starts, or if the person is reading in a moving car. It is caused by a mismatch between visual input and vestibular sensation. For example, the person is reading a book that is stationary in relation to the body, but the vestibular system senses that the car, and thus the body, is moving.
Alternobaric vertigo
Alternobaric vertigo is caused by a pressure difference between the middle ear cavities, usually due to blockage or partial blockage of one eustachian tube, usually when flying or diving underwater. It is most pronounced when the diver is in the vertical position; the spinning is toward the ear with the higher pressure and tends to develop when the pressures differ by 60 cm of water or more.
Decompression sickness
Vertigo is recorded as a symptom of decompression sickness in 5.3% of cases by the U.S. Navy as reported by Powell, 2008 including isobaric decompression sickness.
Decompression sickness can also be caused at a constant ambient pressure when switching between gas mixtures containing different proportions of different inert gases. This is known as isobaric counterdiffusion, and presents a problem for very deep dives. For example, after using a very helium-rich trimix at the deepest part of the dive, a diver will switch to mixtures containing progressively less helium and more oxygen and nitrogen during the ascent. Nitrogen diffuses into tissues 2.65 times slower than helium, but is about 4.5 times more soluble. Switching between gas mixtures that have very different fractions of nitrogen and helium can result in "fast" tissues (those tissues that have a good blood supply) increasing their total inert gas loading. This is often found to provoke inner ear decompression sickness, as the ear seems particularly sensitive to this effect.
Stroke
A stroke (either ischemic or hemorrhagic) involving the posterior fossa is a cause of central vertigo. Risk factors for a stroke as a cause of vertigo include increasing age and known vascular risk factors. Presentation may more often involve headache or neck pain, additionally, those who have had multiple episodes of dizziness in the months leading up to presentation are suggestive of stroke with prodromal TIAs. The HINTS exam as well as imaging studies of the brain (CT, CT angiogram, MRI) are helpful in diagnosis of posterior fossa stroke.
Vertebrobasilar insufficiency
Vertebrobasilar insufficiency, notably Bow Hunter's syndrome, is a rare cause of positional vertigo, especially when vertigo is triggered by rotation of the head.
Management
Definitive treatment depends on the underlying cause of vertigo. People with Ménière's disease have a variety of treatment options to consider when receiving treatment for vertigo and tinnitus including: a low-salt diet and intratympanic injections of the antibiotic gentamicin or surgical measures such as a shunt or ablation of the labyrinth in refractory cases.
Common drug treatment options for vertigo may include the following:
Anticholinergics such as hyoscine hydrobromide (scopolamine)
Anticonvulsants such as topiramate or valproic acid for vestibular migraines
Antihistamines such as betahistine, dimenhydrinate, or meclizine, which may have antiemetic properties
Beta blockers such as metoprolol for vestibular migraine
Corticosteroids such as methylprednisolone for inflammatory conditions such as vestibular neuritis or dexamethasone as a second-line agent for Ménière's disease
All cases of decompression sickness should be treated initially with 100% oxygen until hyperbaric oxygen therapy (100% oxygen delivered in a high-pressure chamber) can be provided. Several treatments may be necessary, and treatment will generally be repeated until either all symptoms resolve, or no further improvement is apparent.
Etymology
Vertigo is from the Latin word, vertō, which means "a whirling or spinning movement".
| Biology and health sciences | Symptoms and signs | Health |
20878009 | https://en.wikipedia.org/wiki/French%20formal%20garden | French formal garden | The French formal garden, also called the , is a style of "landscape" garden based on symmetry and the principle of imposing order on nature. Its epitome is generally considered to be the Gardens of Versailles designed during the 17th century by the landscape architect André Le Nôtre for Louis XIV and widely copied by other European courts.
History
Renaissance influence
The jardin à la française evolved from the French Renaissance garden, a style which was inspired by the Italian Renaissance garden at the beginning of the 16th century. The Italian Renaissance garden, typified by the Boboli Gardens in Florence and the Villa Medici in Fiesole, was characterized by planting beds, or parterres, created in geometric shapes, and laid out symmetrical patterns; the use of fountains and cascades to animate the garden; stairways and ramps to unite different levels of the garden; grottos, labyrinths, and statuary on mythological themes. The gardens were designed to represent harmony and order, the ideals of the Renaissance, and to recall the virtues of Ancient Rome. Additionally, the symmetry of French gardens was a continuation of the Renaissance themes of harmony. French gardens were symmetrical and well manicured to represent order, and this idea of orderliness extended to French society at the time.
Following his campaign in Italy in 1495, where he saw the gardens and castles of Naples, King Charles VIII brought Italian craftsmen and garden designers, such as Pacello da Mercogliano, from Naples and ordered the construction of Italian-style gardens at his residence at the Château d'Amboise and at Château Gaillard, another private résidence in Amboise. His successor Henry II, who had also travelled to Italy and had met Leonardo da Vinci, created an Italian-style garden nearby at the Château de Blois. Beginning in 1528, King Francis I created new gardens at the Château de Fontainebleau, which featured fountains, parterres, a forest of pine trees brought from Provence, and the first artificial grotto in France. The Château de Chenonceau had two gardens in the new style, one created for Diane de Poitiers in 1551, and a second for Catherine de' Medici in 1560.
In 1536 the architect Philibert de l'Orme, upon his return from Rome, created the gardens of the Château d'Anet following the Italian rules of proportion. The carefully prepared harmony of Anet, with its parterres and surfaces of water integrated with sections of greenery, became one of the earliest and most influential examples of the classic French garden. Today, water remains a key garden design in the form of round pools and long ponds.
While the gardens of the French Renaissance were much different in their spirit and appearance than those of the Middle Ages, they were still not integrated with the architecture of the châteaux, and were usually enclosed by walls. In French garden design, the chateau or home was supposed to be the visual focal point. The different parts of the gardens were not harmoniously joined, and they were often placed on difficult sites chosen for terrain easy to defend, rather than for beauty. All this was to change in the middle of the 17th century with the development of the first real garden à la française.
Vaux-le-Vicomte
The first important garden à la française was the Chateau of Vaux-le-Vicomte, created for Nicolas Fouquet, the Superintendent of Finances to Louis XIV, beginning in 1656. Fouquet commissioned Louis Le Vau to design the chateau, Charles Le Brun to design statues for the garden, and André Le Nôtre to create the gardens. It was for the first time that the garden and the chateau were perfectly integrated. A grand perspective of 1500 meters extended from the foot of the chateau to the statue of the Farnese Hercules, and the space was filled with parterres of evergreen shrubs in ornamental patterns, bordered by coloured sand, and the alleys were decorated at regular intervals by statues, basins, fountains, and carefully sculpted topiaries. "The symmetry attained at Vaux achieved a degree of perfection and unity rarely equalled in the art of classic gardens. The chateau is at the center of this strict spatial organization, which symbolizes power and success."
Gardens of Versailles
The Gardens of Versailles, created by André Le Nôtre between 1662 and 1700, were the greatest achievement of the garden à la française. They were the largest gardens in Europe, with an area of 15,000 hectares, and were laid out on an east–west axis followed the course of the sun: the sun rose over the Court of Honor, lit the Marble Court, crossed the Chateau and lit the bedroom of the King, and set at the end of the Grand Canal, reflected in the mirrors of the Hall of Mirrors. In contrast with the grand perspectives, reaching to the horizon, the garden was full of surprises – fountains, small gardens filled with statuary, which provided a more human scale and intimate spaces.
The central symbol of the garden was the sun; the emblem of Louis XIV, illustrated by the statue of Apollo in the central fountain of the garden. "The views and perspectives, to and from the palace, continued to infinity. The king ruled over nature, recreating in the garden not only his domination of his territories, but over the court and his subjects."
Decline
André Le Nôtre died in 1700, but his pupils and his ideas continued to dominate the design of gardens in France through the reign of Louis XV. His nephew, Claude Desgots, created the garden at Château de Bagnolet (Seine-Saint-Denis) for Philippe II, Duke of Orléans (1717) and at Champs (Seine-et-Marne), and another relative, , created gardens for Madame de Pompadour at Crécy (Eure-et-Loir) in 1746 and Bellevue (Hauts-de-Seine) in 1748–50. The major inspiration for gardens continued to be architecture, rather than nature – the architect Ange-Jacques Gabriel designed elements of the gardens at Versailles, Choisy (Val-de-Marne), and Compiègne.
Nonetheless, a few variations in the strict geometry of the garden à la française began to appear. Elaborate parterres of broderies, with their curves and counter-curves, were replaced by parterres of grass bordered with flowerbeds, which were easier to maintain. Circles became ovals, called rotules, with alleys radiating outward in the shape of an 'x', and irregular octagon shapes appeared. Gardens began to follow the natural landscape, rather than moving earth to shape the ground into artificial terraces. Limited colors were available at the time as well. Traditionally, French gardens included blue, pink, white, and mauve.
The middle of the 18th century saw spread in popularity of the new English landscape garden, created by British aristocrats and landowners, and the Chinese style, brought to France by Jesuit priests from the Court of the Emperor of China. These styles rejected symmetry in favor of nature and rustic scenes and brought an end to the reign of the symmetrical garden à la française. In many French parks and estates, the garden closest to the house was kept in the traditional à la française style, but the rest of the park was transformed into the new style, called variously jardin à l'anglaise (the English garden), "anglo-chinois", exotiques, or "pittoresques". This marked the end of the age of the garden à la française and the arrival in France of the jardin paysager, or landscape garden, which was inspired not by architecture but by painting, literature and philosophy.
Theorists and gardeners
Jacques Boyceau, sieur de la Barauderie () the superintendent of royal gardens under Louis XIII, became the first theorist of the new French style. His book, Traité du jardinage selon les raisons de la nature et de l'art. Ensemble divers desseins de parterres, pelouzes, bosquets et autres ornements was published after his death in 1638. Its sixty-one engravings of designs for parterres and bosquets made it a style book for gardens, which influenced the design the Palais du Luxembourg, the Jardin des Tuileries, and the gardens of Saint Germain-en-Laye.
Claude Mollet (ca 1564-shortly before 1649), was the chief gardener of three French kings: Henry IV, Louis XIII, and the young Louis XIV. His father was head gardener at the Château d'Anet, where Italian formal gardening was introduced to France and where Claude apprenticed. His son was André Mollet, who took the French style to the Netherlands, Sweden and England.
André Le Nôtre (1613–1700) was the most important figure in the history of the French garden. The son of the gardener of Louis XIII, he worked on the plans of Vaux-le-Vicomte, before becoming the chief gardener of Louis XIV between 1645 and 1700, and the designer of the Gardens of Versailles, the greatest garden project of the age. The gardens he created became the symbols of French grandeur and rationality, setting the style for European gardens until the arrival of the English landscape park in the 18th century.
Joseph-Antoine Dezallier d'Argenville (1680–1765) wrote Théorie et traité de jardinage, laid out the principles of the garden à la française, and included drawings and designs of gardens and parterres. It was reprinted many times, and was found in the libraries of aristocrats across Europe.
Principles
Jacques Boyceau de La Barauderie wrote in 1638 in his Traité du jardinage, selon les raisons de la nature et de l'art that "the principal reason for the existence of a garden is the esthetic pleasure which it gives to the spectator."
The form of the French garden was largely fixed by the middle of the 17th century. It had the following elements, which became typical of the formal French garden:
a geometric plan using the most recent discoveries of perspective and optics
a terrace overlooking the garden, allowing the visitor to see all at once the entire garden. As the French landscape architect Olivier de Serres wrote in 1600, "It is desirable that the gardens should be seen from above, either from the walls, or from terraces raised above the parterres."
all vegetation is constrained and directed to demonstrate the mastery of man over nature. Trees are planted in straight lines and carefully trimmed, and their tops are trimmed at a set height
the residence serves as the central point of the garden and its central ornament. No trees are planted close to the house; rather, the house is set apart by low parterres and trimmed bushes
a central axis, or perspective, perpendicular to the facade of the house, on the side opposite the front entrance. The axis extends either all the way to the horizon (Versailles) or to piece of statuary or architecture (Vaux-le-Vicomte). The axis faces either South (Vaux-le-Vicomte, Meudon) or east–west (Tuileries, Clagny, Trianon, Sceaux). The principal axis is composed of a lawn, or a basin of water, bordered by trees. The principal axis is crossed by one or more perpendicular perspectives and alleys
the most elaborate parterres, or planting beds, in the shape of squares, ovals, circles or scrolls, are placed in a regular and geometric order close to the house, to complement the architecture and to be seen from above from the reception rooms of the house
the parterres near the residence are filled with broderies, designs created with low boxwood to resemble the patterns of a carpet, and given a polychrome effect by plantings of flowers, or by colored brick, gravel or sand
farther from the house, the broderies are replaced with simpler parterres, filled with grass, and often containing fountains or basins of water. Beyond these, small carefully created groves of trees serve as an intermediary between the formal garden and the masses of trees of the park. "The perfect place for a stroll, these spaces present alleys, stars, circles, theaters of greenery, galleries, spaces for balls and for festivities."
bodies of water (canals, basins) serve as mirrors, doubling the size of the house or the trees
the garden is animated with jeux d'eau and pieces of sculpture, usually on mythological themes, which either underline or punctuate the perspectives, and mark the intersections of the axes, and by moving water in the form of cascades and fountains.
Colours, flowers and trees
Ornamental flowers were relatively rare in French gardens in the 17th century and there was a limited range of colours: blue, pink, white and mauve. Brighter colours (yellow, red, orange) would not arrive until about 1730, because of botanical discoveries from around the world brought to Europe. Bulbs of tulips and other exotic flowers came from Turkey and the Netherlands. An important ornamental feature in Versailles and other gardens was the topiary, a tree or bush carved into geometric or fantastic shapes, which were placed in rows along the main axes of the garden, alternating with statues and vases.
At Versailles flower beds were found only at the Grand Trianon and in parterres on the north side of the palace. Flowers were usually brought from Provence, kept in pots, and changed three or four times a year. Palace records from 1686 show that the palace used 20,050 jonquil bulbs, 23000 cyclamen, and 1700 lily plants.
Most of the trees at Versailles were taken from the forest; they included hornbeam, elm, linden, and beech trees. There were also chestnut trees from Turkey and acacia trees. Large trees were dug up from the forests of Compiègne and Artois and transplanted to Versailles. Many died in transplanting and had to be regularly replaced.
The trees in the park were trimmed both horizontally and flattened at the top, giving them the desired geometric form. Only in the 18th century were they allowed to grow freely.
Parterres de broderie
The parterres de broderie (from the French meaning 'embroidery') is the typical form of French garden design of the Baroque. It is characterised by a symmetrical layout of the flower beds and sheared box hedging to form ornamental patterns known as broderie. Even the arrangement of the flowers is designed to create a harmonious interplay of colours. Frequently found in French Baroque gardens are water gardens, cascades, grottos and statues. Further away from the country house, stately home, chateau or schloss the parterre transitions into the bosquets.
Well known examples are the gardens at the Palace of Versailles in France and the Palace of Augustusburg at Brühl, near Cologne in Germany, which have achieved UNESCO World Heritage status.
As fashions changed, many parterres de broderie of stately homes had to give way in the 19th century to English landscape gardens and have not been reinstated.
Architecture
The designers of the French garden saw their work as a branch of architecture, which simply extended the space of the building to the space outside the walls, and ordered nature according to the rules of geometry, optics and perspective. Gardens were designed like buildings, with a succession of rooms which a visitor could pass through following an established route, hallways, and vestibules with adjoining chambers. They used the language of architecture in their plans; the spaces were referred to as salles, chambres and théâtres of greenery. The "walls" were composed of hedges, and "stairways" of water. On the ground were tapis, or carpets, of grass, brodés, or embroidered, with plants, and the trees were formed into rideaux, or curtains, along the alleys.
Just as architects installed systems of water into the chateaux, they laid out elaborate hydraulic systems to supply the fountains and basins of the garden. Long basins full of water replaced mirrors, and the water from fountains replaced chandeliers. In the bosquet du Marais in the gardens of Versailles, André Le Nôtre placed tables of white and red marble for serving meals. The flowing water in the basins and fountains imitated water pouring into carafes and crystal glasses. The dominant role of architecture in the garden did not change until the 18th century, when the English garden arrived in Europe and the inspiration for gardens began to come not from architecture but from romantic painting.
Theatre
The garden à la française was often used as a setting for plays, spectacles, concerts, and displays of fireworks. In 1664, Louis XIV celebrated a six-day festival in the gardens, with cavalcades, comedies, ballets, and fireworks. Gardens of Versailles included a theatre of water, decorated with fountains and statues of the infancy of the gods (destroyed between 1770 and 1780). Full-size ships were constructed for sailing on the Grand Canal, and the garden had an open-air ballroom surrounded by trees; a water organ, a labyrinth, and a grotto.
Perspective
The architects of the garden à la française did not stop at applying the rules of geometry and perspective to their work. In the first published treatises on gardens, in the 17th century, they devoted chapters to the subject of how to correct or improve perspective, usually to create the illusion of greater distance. This was often done by having alleys become narrower, or having rows of trees that converged, or were trimmed so that they became gradually shorter, as they went farther away from the centre of the garden or from the house. This created the illusion that the perspective was longer and that the garden was larger than it actually was.
Another trick used by French garden designers was the ha-ha (fr: saut de loup). This was a method used to conceal fences which crossed long alleys or perspectives. A deep and wide trench with vertical wall of stone on one side was dug wherever a fence crossed a view, or a fence was placed in bottom of the trench, so that it was invisible to the viewer.
As gardens became more and more ambitious and elaborate through the 17th century, the garden no longer served as a decoration for the chateau. At Chantilly and at Saint-Germain, the chateau became a decorative element of the much larger garden.
Technologies
The appearance of the French garden in the 17th and 18th centuries was a result of the development of several new technologies. The first was géoplastie, the science of moving large amounts of earth. This science had several technological developments. This science had come from the military, following the introduction of cannon and modern siege warfare, when they were required to dig trenches and build walls and earth fortifications quickly. This led to the development of baskets for carrying earth on the back, wheelbarrows, carts and wagons. Andre LeNotre adapted these methods to build the level terraces, and to dig canals and basins on a grand scale.
A second development was in hydrology, bringing water to the gardens for the irrigation of the plants and for use in the many fountains. This development was not fully successful at Versailles, which was on a plateau; even with 221 pumps and a system of canals bringing water from the Seine, and the construction in 1681 of a huge pumping machine, the Machine de Marly, there was still not enough water pressure for all the fountains of Versailles to be turned on at once. Fontainiers were placed along the routes of the King's promenades, and turned on the fountains at each site just before he arrived.
A related development took place in hydroplasie, the art and science of shaping water into different shapes as it came out the fountain. The shape of the water depended upon the force of the water and the shape of the nozzle. New forms created through this art were named tulipe (the tulip), double gerbe (the double sheaf), Girandole(centerpiece) candélabre (candelabra), and corbeille (bouquet), La Boule en l'air (Ball in the air), and L'Evantail (the fan). This art was closely associated with the fireworks of the time, which tried to achieve similar effects with fire instead of water. Both the fountains and fireworks were often accompanied by music, and were designed to show how nature (water and fire) could be shaped by the will of man.
Another important development was in horticulture, in the ability to raise plants from warmer climates in the northern European climate by protecting them inside buildings and bringing them outdoors in pots. The first orangeries were built in France in the 16th century following the introduction of the orange tree after the Italian Wars. The Versailles Orangerie had walls five meters thick, with a double wall that maintains temperatures in winter between . Today it can shelter 1055 trees.
List
Predecessors in the Renaissance Style
Château d'Anet (1536)
Château de Villandry (1536, destroyed in the 19th century and recreated beginning in 1906)
Chateau Fontainebleau (1522–1540)
Château de Chenonceau, gardens of Diane de Poitiers and Catherine de Medici (1559–1570)
Gardens designed by André Le Nôtre
Source:
Vaux-le-Vicomte (1658–1661)
Château de Versailles (1662–1700)
Château de Chantilly (1663–1684)
Château de Fontainebleau (1645–1685)
Château de Saint-Cloud (1664–65)
Gardens of the Tuileries Palace (1664)
Grand Canal of Gardens of Versailles (1668–1669)
Château de Saint-Germain-en-Laye (1669–1673)
Parc de Sceaux (1670)
Château de Dampierre (1673–1783)
Grand Trianon at Versailles (1687–1688)
Château de Clagny (1674–1680)
Château de Meudon
Château de Cordès (1695)
Château de Braine
Château de Pontchartrain
Gardens attributed to André Le Nôtre
Château du Raincy
Château de Courances
Château de Castries
Castle of Racconigi
Later gardens
Château de Lunéville (1710–1724, later an English garden, restored to original design in 2003)
Château de Breteuil (1730–1784)
19th–21st century
Jardin de la Magalone, Marseille, garden by Eduard Andre, 1891.
Nemours Mansion and Gardens – du Pont estate, early 20th century.
Pavillon de Galon in Cucuron, created in 2004
Gardens outside France
Austria
Mirabell Palace in Salzburg
Belvedere Palace in Vienna (designed by Dominique Girard)
Schönbrunn Palace in Vienna (designed by Jean Trehet)
Augarten in Vienna
Parc of Schloss Hof in Engelhartstetten, Lower Austria
Czech Republic
Vrtba Garden, Prague (1720s)
Gardens of the Wallenstein Palace in Prague
England
Blenheim Palace, Oxfordshire (1705–1724)
The Parterre, Waddesdon Manor, Buckinghamshire (1870s)
Germany
Schwetzingen Palace in Schwetzingen, Baden-Württemberg
Weikersheim Castle in Weikersheim, Baden-Württemberg
Ludwigsburg Palace near Stuttgart, Baden-Württemberg
Gardens of the Würzburg Residence in Würzburg, Bavaria
Schleissheim Palace in Munich, Bavaria
Nymphenburg Palace in Munich, Bavaria
Karlsaue, Kassel, Hesse (built until 1785)
French Garden, Celle in Celle, Lower Saxony
Herrenhausen Gardens, Hanover, Lower Saxony (1676–1680)
Augustusburg and Falkenlust Palaces, Brühl in Brühl (Rhineland), North Rhine-Westphalia
French garden of Schloss Benrath in Düsseldorf, North Rhine-Westphalia
Italy
Royal Palace of Caserta near Napoli
Palazzina di caccia of Stupinigi Palace, Piedmont
Racconigi Palace, Piedmont (1755)
Netherlands
Het Loo Palace in Apeldoorn, Gelderland
Poland
Parc of Nieborów Palace, Łódź Voivodeship (designed by Tylman van Gameren)
Branicki Palace, Białystok, Podlaskie Voivodeship (1737–1771)
Russia
Peterhof Gardens, St. Petersburg (1714–1725)
Summer Garden, St. Petersburg (1712–1725)
Tsarskoe Selo Old Garden in Pushkin (1717–1720)
Kuskovo Estate, Moscow (1750–1780)
Oranienbaum Palace and Garden, west of St. Petersburg
Spain
Royal Palace of La Granja de San Ildefonso in San Ildefonso, Segovia
Sweden
Drottningholm Palace gardens outside Stockholm
| Technology | Buildings and infrastructure | null |
10660257 | https://en.wikipedia.org/wiki/Pendant%20group | Pendant group | In IUPAC nomenclature of chemistry, a pendant group (sometimes spelled pendent) or side group is a group of atoms attached to a backbone chain of a long molecule, usually a polymer. Pendant groups are different from pendant chains, as they are neither oligomeric nor polymeric.
For example, the phenyl groups are the pendant groups on a polystyrene chain.
Large, bulky pendant groups such as adamantyl usually raise the glass transition temperature () of a polymer by preventing the chains from sliding past each other easily. Short alkyl pendant groups may lower the by a lubricant effect.
| Physical sciences | Concepts_2 | Chemistry |
1881328 | https://en.wikipedia.org/wiki/Juniperus%20oxycedrus | Juniperus oxycedrus | Juniperus oxycedrus, vernacularly called Cade, cade juniper, prickly juniper, prickly cedar, or sharp cedar, is a species of juniper, native across the Mediterranean region, growing on a variety of rocky sites from sea level. The specific epithet oxycedrus means "sharp cedar" and this species may have been the original cedar or cedrus of the ancient Greeks.
Description
Juniperus oxycedrus is very variable in shape, forming a spreading shrub tall to a small erect tree tall. It has needle-like leaves in whorls of three; the leaves are green, long and broad, with a double white stomatal band (split by a green midrib) on the inner surface. It is usually dioecious, with separate male and female plants. The seed cones are berry-like, green ripening in 18 months to orange-red with a variable pink waxy coating; they are spherical, diameter, and have three or six fused scales in 1–2 whorls, three of the scales with a single seed. The seeds are dispersed when birds eat the cones, digesting the fleshy scales and passing the hard seeds in their droppings. The pollen cones are yellow, long, and fall soon after shedding their pollen in late winter or early spring.
Subspecies
As to be expected from the wide range, J. oxycedrus is very variable, and multiple subspecies have been recognised. However, multiple studies have found the subspecies not to be closely related to one another, resulting in the recognition of multiple species:
Juniperus oxycedrus L. – Western prickly juniper. Southwest Europe, in eastern Portugal and Spain east to southern France, northwest Italy, Corsica, and Sardinia, and northwest Africa from Morocco east to Tunisia. Leaves long (), narrow-based; cones smooth.
Juniperus navicularis Gand. (syn. J. oxycedrus subsp. transtagana) – Portuguese prickly juniper. Coastal southwest Portugal. Leaves short (); cones smooth.
Juniperus deltoides R.P.Adams – Eastern prickly juniper. Central Italy east to Iran and Israel. Leaves long (), broad-based; cones with raised scale edges.
Juniperus macrocarpa (syn. J. oxycedrus subsp. macrocarpa) – large-fruited juniper. Mediterranean coastal sands. Broader leaves ( wide), and larger cones ( wide).
An additional variety or subspecies J. oxycedrus var. badia H.Gay (syn. J. oxycedrus subsp. badia (H.Gay) Debeaux) is distinguished on the basis of larger cones ( diameter), tinged purple when mature; it is described from northern Algeria, and also reported from Portugal and Spain.
Other close relatives of J. oxycedrus include Juniperus brevifolia on the Azores, Juniperus cedrus on the Canary Islands and Juniperus formosana in eastern Asia.
Uses
Cade oil is the essential oil obtained through destructive distillation of the wood of this shrub. It is a dark, aromatic oil with a strong smoky smell which is used in some cosmetics and (traditional) skin treatment drugs, as well as incense. Cade oil has, on rare occasions, caused severe allergic reactions in infants.
| Biology and health sciences | Cupressaceae | Plants |
1881635 | https://en.wikipedia.org/wiki/Chironex%20fleckeri | Chironex fleckeri | Chironex fleckeri, commonly known as the Australian box jelly, and nicknamed the sea wasp, is a species of extremely venomous box jellyfish found in coastal waters from northern Australia and New Guinea to Indonesia, Cambodia, Malaysia and Singapore, the Philippines and Vietnam. It has been described as "the most lethal jellyfish in the world", with at least 64 known deaths in Australia from 1884 to 2021.
Notorious for its sting, C. fleckeri has tentacles up to long covered with millions of cnidocytes which, on contact, release microscopic darts delivering an extremely powerful venom. Being stung commonly results in excruciating pain, and if the sting area is significant, an untreated victim may die in two to five minutes. The amount of venom in one animal is said to be enough to kill 60 adult humans.
Taxonomy
Chironex fleckeri was named after North Queensland toxicologist and radiologist Doctor Hugo Flecker. On January 20, 1955, when a 5-year-old boy died after being stung in shallow water at Cardwell, North Queensland, Flecker found three types of jellyfish. One was an unidentified box-shaped jellyfish with groups of tentacles arising from each corner. Flecker sent it to Dr. Ronald Southcott in Adelaide, and on December 29, 1955, Southcott published his article introducing it as a new genus and species of lethal box jellyfish. He named it Chironex fleckeri: the genus name a compound of the centaur Chiron in Greek mythology (whose name is probably related to Ancient Greek χείρ meaning "hand") and Latin nex meaning "murder", and the specific epithet "fleckeri" in honour of its discoverer. The name "hand of death" refers to the four appendages of C. fleckeri appearing as hands.
Description
Chironex fleckeri is the largest of the cubozoans (collectively called box jellyfish), many of which may carry similarly toxic venom. Its bell usually reaches about in diameter but can grow up to . From each of the four corners of the bell trails a cluster of 15 tentacles. The pale blue bell has faint markings; viewed from certain angles, it bears a somewhat eerie resemblance to a human head or skull. Since it is virtually transparent, the creature is nearly impossible to see in its habitat, posing significant danger to swimmers.
When the jellyfish are swimming, the tentacles contract so they are about long and about in diameter; when they are hunting, the tentacles are thinner and extend to about long. The tentacles are covered with a high concentration of stinging cells called cnidocytes, which are activated by pressure and a chemical trigger; they react to proteinous chemicals. Box jellyfish are day hunters; at night they are seen resting on the ocean floor.
In common with other box jellyfish, C. fleckeri has four eye-clusters with 24 eyes. Some of these eyes seem capable of forming images, but whether they exhibit any object recognition or object tracking is debated; it is also unknown how they process information from their sense of touch and eye-like light-detecting structures due to their lack of a central nervous system. They are attracted to light of different colors (white, red, orange, yellow, green and blue), but blue light seems to elicit a feeding behavior, as it slows down their pulsation rate and makes them stream out their tentacles. Black objects, on the other hand, cause them to move away.
Chironex fleckeri lives on a diet of prawns and small fish, as well as crabs and other pelagic invertebrates. Their only known predators are green sea turtles and leatherback turtles, whose thick skin is impenetrable to the cnidocytes of the jellyfish, among other pelagic predators.
Distribution and habitat
The medusa is pelagic and has been documented from coastal waters of Australia and New Guinea north to the Philippines, Malaysia, Singapore and Vietnam. In Australia, it is known from the northern coasts from Exmouth to Agnes Water, but its full distribution outside Australia has not been properly identified. To further confuse, the closely related and also dangerously venomous Chironex yamaguchii was first described from Japan in 2009. This species has also been documented from the Philippines, meaning the non-Australian records of C. fleckeri need to be rechecked. Breeding occurs in lower levels of rivers and mangrove channels.
Sting
Chironex fleckeri is best known for its extremely powerful "sting". The sting can produce excruciating pain accompanied by an intense burning sensation, like being branded with a red hot iron. Stings may also result in white welts and lines that may be accompanied by blistering. In some cases, the sting can cause permanent damage or death to the skin and result in scars. Fatalities are most often caused by specimens of C. fleckeri that are larger than . If left untreated, large amounts of venom injection can cause fatality in 5 minutes.
In Australia, C. fleckeri has caused at least 64 deaths since the first report in 1883, but most encounters appear to result only in mild envenomation. Among 225 analyzed C. fleckeri stings in Australia's Top End from 1991 to 2004, only 8% required hospital admission, 5% received antivenom and there was a single fatality (a 3-year-old child). 26% experienced severe pain, while it was moderate to none in the remaining. Most deaths in recent decades have been children, as their smaller body mass puts them at a higher risk of fatal envenomation. When people do die, it is usually caused by a cardiac arrest occurring within minutes of the sting. It takes approximately of tentacle to deliver the fatal dose.
The venom causes cells to become porous enough to allow potassium leakage, causing hyperkalemia, which can lead to cardiovascular collapse and death as quickly as within two to five minutes with an of 0.04 mg/kg. It was postulated that a zinc compound may be developed as an antidote. Occasionally, swimmers who get stung will undergo cardiac arrest or drown before they can even get back to the shore or boat.
In many cases, there will be a reaction that takes place days after the initial sting if the victim survives. This extremely itchy rash can last weeks after the initial sting. If the skin in the affected area is intact, certain creams and antihistamines may help to alleviate the symptoms.
Chironex fleckeri and other jellyfish, including the Irukandji (Carukia barnesi), are abundant in the waters of northern Australia during the warmer months of the year. They are believed to drift into estuaries to breed. Signs like the one pictured are erected along the coast of North Queensland to warn people of such, and few people swim during this period. Some people still do, however, putting themselves at great risk. At popular swimming spots, net enclosures are placed out in the water wherein people can swim but jellyfish cannot get in, keeping swimmers safe.
History of sting treatment
Until 2005, treatment involved using pressure immobilisation bandages, with the aim of preventing distribution of the venom through the lymph and blood circulatory systems. This treatment is no longer recommended by health authorities, due to research which showed that using bandages to achieve tissue compression provoked nematocyst discharge.
The application of vinegar is recommended treatment because vinegar (4–6% acetic acid) permanently deactivates undischarged nematocysts, preventing them from opening and releasing venom. A 2014 study demonstrated in vitro that while vinegar deactivates unfired nematocysts, there was also an increase in venom concentration in the solution, possibly by causing already-fired nematocysts (which still contain some venom) to release what remained. However, this study has been criticized on several methodological grounds, including that the experiment was done using a model membrane that is much different from (and more simple than) human skin. Also, the researchers did not determine whether the increase in venom concentration was caused by already-discharged nematocysts releasing more venom, or if the venom that was released initially had simply leaked back out through the membrane, thus confounding the concentration measurement. Despite these concerns, diluted acetic acid is still the recommended treatment. In 2019, the first antidote for Australian box jellyfish sting was discovered in Australia.
Reproduction
Chironex fleckeri is capable of both sexual and asexual reproduction and are oviparous.
Sexual Reproduction
A fully grown and sexually mature Chironex fleckeri medusa will begin trying to find a mate in the spring season. Usually, the jellies will move from their usual habitat to a freshwater river for this hunt. If a mate is located, sperm and eggs are expelled into the water by the male and female respectively to result in fertilization. The gametes produced by the spawning process will go on to become planulae and eventually small sea wasp polyps. They will use their two tentacles to hide away from predators by hooking onto a hidden surface and feeding on plankton. The parent organisms die shortly after reproduction occurs.
Asexual Reproduction (Budding)
Polyps created through spawning will begin to form entirely new clones of themselves through a process called budding. The original polyp along with the new ones will eventually become small medusa through metamorphosis and return to the ocean until they are mature. The process will then repeat once they can sexually reproduce and go back to a river to find a mate.
| Biology and health sciences | Cnidarians | Animals |
1881649 | https://en.wikipedia.org/wiki/Ulmus%20laevis | Ulmus laevis | Ulmus laevis Pall., variously known as the European white elm, fluttering elm, spreading elm, stately elm and, in the United States, the Russian elm, is a large deciduous tree native to Europe, from France northeast to southern Finland, east beyond the Urals into Kyrgyzstan and Kazakhstan, and southeast to Bulgaria and the Crimea; there are also disjunct populations in the Caucasus and Spain, the latter now considered a relict population rather than an introduction by man, and possibly the origin of the European population. U. laevis is rare in the UK, although its random distribution, together with the absence of any record of its introduction, has led at least one British authority to consider it native. NB: The epithet 'white' elm commonly used by British foresters alluded to the timber of the wych elm.
The species was first identified, as Ulmus laevis, by Pallas, in his Flora Rossica published in 1784. The tree is allogamous and is most closely related to the American elm U. americana.
Endemic to alluvial forest, U. laevis is rarely encountered at elevations above 400 m. Most commonly found along rivers such as the Volga and Danube, it is one of very few elms tolerant of prolonged waterlogged, anoxic ground conditions. The species is threatened by habitat destruction and disturbance in some countries, notably Spain. Flood control schemes are particularly harmful, as seed dispersion is reliant on floods, while abstraction from aquifers lowering ground water levels has compromised the development of the trees.
Although not possessed of an innate genetic resistance to Dutch elm disease, the species is rarely infected in western Europe.
Description
Ulmus laevis is similar in stature to the wych elm, if rather less symmetric, with a looser, untidy, branch structure and less neatly rounded crown. The tree typically reaches a height and breadth of > 30 m, with a trunk < 2 m d.b.h. The extensive shallow root system ultimately forms distinctive high buttresses around the base of the trunk. The bark is smooth at first, then in early maturity breaks into thin grey scales, which separate with age into a network of grey-brown scales and reddish-brown underbark, and finally is deeply fissured in old age like other elms. The leaves are deciduous, alternate, simple ovate with a markedly asymmetric base, < 10 cm long and < 7 cm broad, comparatively thin, often almost papery in texture and very translucent, smooth above with a downy underside. Significantly, the leaf veins do not divide from the central vein to the leaf margin. The leaves are shed earlier in autumn than other species of European elm.
The tree is most reliably distinguished from other European elms by its long flower stems, averaging 20 mm. Moreover, the apetalous wind-pollinated flowers are distinctively cream-coloured, appearing before the leaves in early spring in clusters of 15-30; they are 3–4 mm across. The fruit is a winged samara < 15 mm long by 10 mm broad with a ciliate margin, the single round 5 mm seed maturing in late spring. The seeds have a generally high rate of germination, 45–60% for Serbian trees examined by Stilinović.
Although the species is protandrous, levels of self-pollination can be high The tree can grow very rapidly; where planted in persistently moist soil, trunk width of 13-year-old trees increased by 4 cm per annum at breast height (d.b.h.). The species differs from its closest relative, the American elm, mainly in the irregular crown structure and frequent epicormic shoots, features which also give the tree a distinctive winter silhouette. The American elm also has less acute leaf buds, longer petioles, narrower leaves, and a deeper apical notch in the samara which reaches the seed.
Pests and diseases
Like other European elms, natural populations of the European white elm have little innate resistance to Dutch elm disease. In a study in France, losses to DED amounted to 28% over a 10 year period. However, research by Irstea has isolated clones able to survive injection with the causal fungus, initially losing < 70% of their foliage, but regenerating strongly the following year.
The tree is not favoured by the vector bark beetles, which colonize it only when there are no other elm alternatives available, an uncommon situation in western Europe. Indeed, in a study of elm in Flanders, not one example of U. laevis was found to be afflicted by Dutch elm disease. Research in Spain has indicated that it is the presence of an antifeedant triterpene, alnulin, at a concentration of 200 μg/g {dried bark} which renders the tree unattractive to the beetles. Ergo: the tree's decline in western Europe has been chiefly owing to woodland clearance in river valleys, and river management systems eliminating flooding, not disease. However, in 2020, it was noted by the Dutch forestry commission that many laevis, but only in Zeeland, were succumbing to Dutch elm disease for reasons unclear.
It was noted by Jouin at Metz, and a century later by Mittempergher and Santini in Italy, that U. laevis had a very low susceptibility to the elm leaf beetle Xanthogaleruca luteola. Research in Germany has established that the tree is also eschewed by the Zig Zag sawfly Aproceros leucopoda.
Elwes observed that trees planted at Ugbrooke in Devon were infested with Cacopsylla ulmi, which he had never found on any other elm in Britain, an affliction confirmed many years later by Richens, who discovered the specimens of U. laevis grown at Kew were the only elms in the Gardens afflicted by the louse, and the aphid Tinocallis platani.
The species has a slight to moderate susceptibility to elm yellows.
Cultivation
U. laevis is essentially a riparian tree, able to withstand over 100 days of continual flooding, although it is intolerant of saline conditions Spanish trees were found to be calcifuge, preferring slightly acid, siliceous soils, and also drought-intolerant, their xylem vessels prone to drought-stress cavitation. In England, the tree failed to prosper in chalk stream valleys, where the soil was predominantly black peat, named 'Adventurers' for the Adventurers' Land SSSI in Cambridgeshire, owing to dehydration in summer. Trees planted in dry ground are notoriously short-lived.
U. laevis is comparatively weak-wooded, much more so than field elm Ulmus minor, and thus an inappropriate choice for exposed locations. In trials in southern England by Butterfly Conservation, young trees of <5 m height were badly damaged by wind gusts of 40 knots (75 km/h) in midsummer storms.
The species was never widely introduced to the United States, but is represented at several arboreta. Ulmus effusa, supplied by the Späth nursery of Berlin, was planted at the Dominion Arboretum, Ottawa, Canada, in 1896, as U. pedunculata. In the Far East, the tree has been planted in Xinjiang province and elsewhere in northern China; planting in Tongliao City is known to have been particularly successful. White elm is also known to have been introduced to Australia.
Since the beginning of the 21st century, the tree has enjoyed a small renaissance in England. A popular larval host plant of the white-letter hairstreak Satyrium w-album butterfly across Europe, the elm is now being planted by Butterfly Conservation and other groups to restore local populations decimated by the effects of Dutch elm disease on native or archaeophytic elms. The Cheshire Wildlife Trust, for example, planted numerous white elms on its reserves in the former Vale Royal district of the county.
Introduction to the UK and Ireland
U. laevis is probably not native to the United Kingdom despite its random occurrence in the countryside, although the date and circumstances of its introduction have not survived. The earliest published references to the tree (as U. effusa, citing Willdenow) were in Sibthorp's Flora Oxoniensis (1794), and (as U. effusa Willd. but without description) in Miller's posthumously revised Gardener's and Botanist's Dictionary (1807). The first specimen to be reported in cultivation, in 1838, was at Whiteknights Park, Reading, which featured an elm grove; the tree measured in height, suggesting it had been planted at the end of the 18th century. However, the authenticity of the Whiteknights tree is a matter of contention; it flowered but did not set fertile seed, which suggested to Loudon that it might be U. campestris (U. minor 'Atinia'), or, on account of it not producing suckers, possibly U. montana (:U. glabra). Moreover, Whiteknights was supplied by the Lee and Kennedy nursery of Hammersmith, which is not known to have stocked U. laevis. A tree at Syon Park identified by Elwes & Henry as U. laevis was later considered by Bean as more closely resembling U. americana by dint of its symmetrical branch arching. The species was not reported from the wild until 1943, with the discovery of a tree in a Surrey hedgerow.
It is possible the tree's distribution was associated with Capability Brown (1716–1783), known to have favoured U. laevis, which he listed among his preferred "native" (sic) trees. This could explain the existence of the seven old specimens discovered by Elwes in 1908 on Mount Pleasant within Ugbrooke Park, Devon, designed by Brown in 1761. Ugbrooke is four miles from Mamhead Park, which had earlier been planted with numerous exotic trees, notably holm oak, collected by its owner, merchant Thomas Ball ( d. 1749) during his commercial travels in Europe. Ball's introductions were known to have been marketed by his head gardener William Lucombe, who in 1720 founded the first commercial nursery in the south-west at Exeter, though an account of trees growing at Mamhead by Pince (grandson of Lucombe) in 1835 makes no mention of U. laevis nor of any other elms. None of Lucombe's early catalogues are known to survive, and thus the introduction of U. laevis through south Devon cannot be confirmed. However, the tree does not feature in any of the surviving arboreta accessions lists, or catalogues of the larger, nationally famous, nurseries of the day, and its earliest-known mention in commerce remains in the south-west, in the catalogue of the Ford & Please nursery (as U. pedunculata) at Exeter circa 1836. James Main mentions the tree as 'a native of Hungary' and in 1838 only to be met in 'ornamental plantations', but by 1846 was 'becoming available in (UK) nurseries'.
U. laevis, obtained from the Späth nursery of Berlin as U. effusa, was planted in Kew Gardens (1895), in the Ryston Hall arboretum, Norfolk (1914), and, re-propagated, in Cambridge University Botanic Garden (1909). Evidently the tree did not gain in popularity, and was overlooked or ignored by most authors of popular guides to trees in Britain during the 20th century, notably Mabey in his Flora Britannica. The tree is also omitted from Keble-Martin's comprehensive Flora of Devon.
It is not known whether U. laevis was introduced to Scotland before the early 20th century. Two of the three specimens supplied by the Späth nursery, Berlin, to the Royal Botanic Garden Edinburgh in 1902 as U. effusa may survive in Edinburgh, as it was the practice of the garden to distribute trees about the city; the third specimen was in the garden itself. Other examples can be found in the city, notably in Fettes College grounds opposite Inverleith Allotments, and at the entrance to North Merchiston Cemetery. In Wales, two mature trees with numerous seedlings occur in a small wood at Rhydyfelin near Aberystwyth, while another grows at Llandegfan, Anglesey. In Ireland, the tree is represented by a line of four at the Old Rectory, Kells Road, Ardee, County Louth (, girth , October 2009), and in the Channel Islands, by a clump near the well at La Seigneurie (Le Manoir), Sark.
Notable trees
The two largest known trees in Europe are at Gülitz in Germany (3.1 m d.b.h.), and at Komorów in Poland (2.96 m d.b.h. in 2011), known as the Witcher. Other veterans survive at Casteau, Belgium (bole-girth 5.15 m), in Rahnsdorf near Berlin (bole-girth 4.5 m) and in Ritvala, Finland (bole-girth 4.49 m).
A lane of Ulmus laevis is found at Eibergen, Netherlands (see Gallery below), while a large, mature specimen is found within the Alhambra, Granada.
Ulmus laevis has very occasionally been planted as an ornamental tree in the UK, and even more randomly in countryside hedgerows. The UK Champion is at Ferry Farm, on the banks of the Tamar at Harewood, Cornwall (27 m high, 1.8 m d.b.h. in 1997). Other examples are few and far between though sometimes of considerable age, surviving amid diseased native elm in Cornwall at Torpoint, and Pencalenick (21 m high, d.b.h. 1.75 m), and near Over Wallop in Hampshire (16 m high, d.b.h. 1.3 m 2016) The largest-known aggregation in England is the ring of 50 trees planted circa 1950 within a ring of common lime around a former ammunition dump on the elevated chalk of Salisbury Plain at Hexagon Wood, Larkhill, about north of Stonehenge.
In the United States, a tree of 31.4 m (103 ft) in height (2015) grows at 3331 NE Hancock Street in Portland, Oregon; its age is not known.
Uses
In Finland, young European white elm trees were traditionally grown for the raw material of shaft bows. It's leathery bark was also used in tough bindings.
The density of the timber is significantly lower than that of other European elms. However, owing to its rapid growth, tolerance of soil compaction, air pollution and de-icing salts, the tree has long been used for amenity planting in towns and along roadsides.
Propagation
U. laevis is easily grown from seed sown on ordinary compost and kept well-watered. However, viability can vary greatly from year to year, while the seed is remarkably short-lived. Germination should occur within one week even without heat, the best seedlings attaining as much as half a metre in their first year. Softwood cuttings taken in June is also a reliable method; the cuttings strike very quickly, well within a fortnight, rapidly producing a dense matrix of roots.
Subspecies and varieties
Several putative varieties have been identified. A var. celtidea from Ukraine was reported by Rogowicz in the middle of the 19th century, but no examples are known to survive. Another, var. parvifolia, has been reported from Serbia. A third, var. simplicidens, is very rare; the only example known to survive is at the National Botanic Garden of Latvia in Salaspils. Kew had a grafted var. glabra in the early 20th century (provenance unknown), a clone of which is present at Wakehurst Place.
Cultivars
Compared with the other European species of elm, U. laevis has received scant horticultural attention, there being only eight recorded cultivars:
In Russia other ornamental forms are recognized: f. argentovariegata, f. rubra, and f. tiliifolia. A pyramidal form was reported in 1888 from the Fredericksfelde cemetery in Berlin by Bolle. A line of similar monopodial trees grows (2019) on the island in the Lot at Entraygues, France.
Hybrids
U. laevis does not hybridize naturally, in common with the American elm (U. americana) to which it is closely related. However, in experiments at the Arnold Arboretum, it was successfully crossed with U. thomasii and U. pumila; no such crosses have ever been released to commerce.
Accessions
Europe
Arboretum de La Petite Loiterie, Monthodon, France. No details available
Arboretum Freiburg-Günterstal, Germany, no details available
Brighton & Hove City Council, UK, NCCPG Elm Collection. Ten trees at Hove Recreation Ground, Hove.
Copenhagen University Botanic Garden, Denmark. No details available.
ELTE Botanic Garden, Budapest, Hungary. Acc. nos. 1998-0718, 1998-0719.
Grange Farm Arboretum, Sutton St. James, Spalding, Lincolnshire, UK. Acc. no. 502.
Great Fontley Butterfly Conservation Elm Trials plantation, UK. Two planted 2003, grown from cuttings of specimen at RBG Wakehurst Place.
Hortus Botanicus Nationalis, Salaspils, Latvia. Acc. nos. 18136, 18140.
Linnaean Gardens of Uppsala, Sweden. Acc. no. 1930-1014.
Royal Botanic Garden Edinburgh, UK. Acc. no. 20070643, from seed wild collected in Val d'Allier, France.
Royal Botanic Gardens Kew, UK. Acc. nos. 1969-17302, 1973-11712.
Royal Botanic Gardens, Wakehurst Place, UK. Acc. no. 1973-21048.
Sir Harold Hillier Gardens, Romsey, Hampshire. UK. Acc. no. 2016.0385
Tallinn Botanic Garden, Estonia. No accession details available.
Thenford House arboretum, Northamptonshire, UK. No details available.
'The Leys', University Parks, Oxford, UK. Acc. no. 02678.
Westonbirt Arboretum, UK. Tetbury, Glos., UK. Acc. no. 1995/322
Wijdemeren City Council, Netherlands. Elm Collection. Planted 1990 Tjalk, Loosdrecht; 2007 Hinderdam, Nederhorst den Berg; elm lane De Kwakel, Kortenhoef in 2009.
North America
Arnold Arboretum, US. Acc. nos. 17910, 637-79, 6951, 753-80.
Brenton Arboretum, Dallas Center, Iowa, US. No details available.
Brooklyn Botanic Garden, New York City, US. Acc. no. X02589.
Dominion Arboretum, Canada. No details available
Longwood Gardens, US. Acc. nos. 1964-0568, 1964-1119.
Morton Arboretum, Illinois, US. Acc. nos. 1302-27, 446-48, 492-64, 27-98.
Nurseries
Arboretum Waasland, Nieuwkerken-Waas, Belgium
Boomkwekerij Oirschot, Oirschot, Netherlands
Landford Trees, Salisbury, UK.
Lorenz von Ehren, Hamburg, Germany.
Noordplant, Glimmen, Netherlands
Pan-Global Plants, Frampton-on-Severn, Gloucestershire, UK
UmbraFlor, Spello, Italy
Van Den Berk (UK) Ltd., , London, UK
| Biology and health sciences | Rosales | Plants |
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