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1882227 | https://en.wikipedia.org/wiki/Anabaena | Anabaena | Anabaena is a genus of filamentous cyanobacteria that exist as plankton. They are known for nitrogen-fixing abilities, and they form symbiotic relationships with certain plants, such as the mosquito fern. They are one of four genera of cyanobacteria that produce neurotoxins, which are harmful to local wildlife, as well as farm animals and pets. Production of these neurotoxins is assumed to be an input into its symbiotic relationships, protecting the plant from grazing pressure.
A DNA sequencing project was undertaken in 1999, which mapped the complete genome of Anabaena, which is 7.2 million base pairs long. The study focused on heterocysts, which convert nitrogen into ammonia. Certain species of Anabaena have been used on rice paddy fields, proving to be an effective natural fertilizer.
Nitrogen fixation by Anabaena
Under nitrogen-limiting conditions, vegetative cells differentiate into heterocysts at semiregular intervals along the filaments. Heterocyst cells are terminally specialized for nitrogen fixation. The interior of these cells is micro-oxic as a result of increased respiration, inactivation of O2-producing photosystem (PS) II, and formation of a thickened envelope outside of the cell wall. Nitrogenase, sequestered within these cells, transforms dinitrogen into ammonia at the expense of ATP and reductant—both generated by carbohydrate metabolism, a process supplemented, in the light, by the activity of PS I. Carbohydrate, probably in the form of glucose, is synthesized in vegetative cells and moves into heterocysts. In return, nitrogen fixed in heterocysts moves into the vegetative cells, at least in part in the form of amino acids.
The fern Azolla forms a symbiotic relationship with the cyanobacterium Anabaena azollae, which fixes atmospheric nitrogen, giving the plant access to this essential nutrient. This has led to the plant being dubbed a "super-plant", as it can readily colonise areas of freshwater, and grow at great speed - doubling its biomass in as little as 1.9 days. The typical limiting factor on its growth is phosphorus, abundance of which, due to chemical runoff, often leads to Azolla blooms. Unlike other known plants, the symbiotic microorganism is transferred directly from one generation to the next. This has made Anabaena azollae completely dependent on its host, as several of its genes are either lost or have been transferred to the nucleus in Azolla's cells.
Primitive vision pigments studied in Anabaena
Anabaena is used as a model organism to study simple vision. The process in which light changes the shape of molecules in the retina, thereby driving the cellular reactions and signals that cause vision in vertebrates, is studied in Anabaena. Anabaena sensory rhodopsin, a specific light-sensitive membrane protein, is central to this research.
DNA repair
Double strand breaks (DSBs) are a type of DNA damage that can be repaired by homologous recombination. This enzymatic repair process occurs in several enzymatic steps including an early step catalyzed by RecN protein. A study of the dynamics of RecN in DSB repair in Anabaena indicated differential regulation of DSB repair so that it is active in vegetative cells but absent in mature heterocysts that are terminal cells.
| Biology and health sciences | Gram-negative bacteria | Plants |
1883867 | https://en.wikipedia.org/wiki/Trim%20tab | Trim tab | Trim tabs are small surfaces connected to the trailing edge of a larger control surface on a boat or aircraft, used to control the trim of the controls, i.e. to counteract hydro- or aerodynamic forces and stabilise the boat or aircraft in a particular desired attitude without the need for the operator to constantly apply a control force. This is done by adjusting the angle of the tab relative to the larger surface.
Changing the setting of a trim tab adjusts the neutral or resting position of a control surface (such as an elevator or rudder). As the desired position of a control surface changes (corresponding mainly to different speeds), an adjustable trim tab will allow the operator to reduce the manual force required to maintain that position—to zero, if desired. Thus the trim tab acts as a servo tab. Because the center of pressure of the trim tab is farther away from the axis of rotation of the control surface than the center of pressure of the control surface, the moment generated by the tab can match the moment generated by the control surface. The position of the control surface on its axis will change until the torques from the control surface and the trim surface balance each other.
On boats
Planing boats or boats that operate at speeds close to planing will often have trim tabs on their engine lower unit or attached to the transom. Adjusting them up or down alters the pitch attitude of the boat while under way, variously balancing speed, weight distribution, and sea conditions.
On aircraft
Trim tabs are variously integrated into the rudder, elevators, and ailerons of a fixed-wing aircraft. As such, they are elements of an aircraft's system for allowing the pilot to control and maintain airspeed with a minimum of inputs and mental concentration. Many newer aircraft, especially jet aircraft, have electric trim controls.
Elevator trim frees the pilot from exerting constant force on the pitch controls, by adjusting trim control (often in the form of a vertical wheel) to cancel out control forces for a given airspeed and weight distribution. Typically, when this wheel is rotated up (or lever raised) the aircraft's nose pitches down; rotating it down (or depressing the lever) lowers the tail and raises the nose.
Many airplanes also have rudder and/or aileron trim systems. On some, the rudder trim tab is hinged and adjustable during flight; on others it is only adjustable on the ground (to lessen the need for the pilot to push the rudder pedal constantly to overcome the left-turning tendencies of many prop-driven aircraft).
Most fixed-wing aircraft have a trim tab on the elevator. However, alternative means of controlling the speed and attitude of the aircraft are sometimes used, including:
a spring included in the control system that can be adjusted by the pilot
in the case of the elevator, an all-moving horizontal stabilizer, called a stabilator, the position of which can be adjusted in flight by a servo tab or an anti-servo tab.
On some aircraft (e.g. Concorde, McDonnell Douglas MD-11), fuel may be shifted to tanks in the tail during cruise to reposition the center of gravity in order to reduce trim drag. Maintaining the center of gravity near the aft-most limit for cruise improves cruise efficiency.
When a servo tab is employed, it is moved into the slipstream opposite to the control surface's desired deflection. For example, in order to trim an elevator to hold the nose down, the elevator's trim tab will actually rise up into the slipstream. The increased pressure on top of the trim tab surface caused by raising it will then deflect the entire elevator slab down slightly, causing the tail to rise and the aircraft's nose to move down. In the case of an aircraft where deployment of high-lift devices (flaps) would significantly alter the longitudinal trim, a supplementary trim tab is arranged to simultaneously deploy with the flaps so that pitch attitude is not markedly changed.
The use of trim tabs significantly reduces pilots' workload during continuous maneuvers (e.g. sustained climb to altitude after takeoff or descent prior to landing), allowing them to focus their attention on other tasks such as traffic avoidance or communication with air traffic control.
Both elevator trim and pitch trim affect the small trimming part of the elevator on jet airliners. The former is supposed to be set in a certain position for a longer time, while the pitch trim (controlled with the landing pilot's thumb on the yoke or joystick, and thereby easy to maneuver) is used all the time after the flying pilot has disabled the autopilot, especially after each time the flaps are lowered or at every change in the airspeed, at the descent, approach and final. Elevator trim is most used for controlling the attitude at cruising by the autopilot.
Beyond reducing pilot workload, proper trim also increases fuel efficiency by reducing drag. For example, propeller aircraft have a tendency to yaw when operating at high power, for instance when climbing; this increases parasitic drag because the craft is not flying straight into the apparent wind. In such circumstances, the use of an adjustable rudder trim tab can reduce yaw.
Military
On military aircraft during wartime, trim tabs often served as unintentional backup control systems for aircraft with damaged controls. Since trim tabs are usually controlled by their own dedicated system of control cables, rods, and/or hydraulic lines, aircraft that had suffered loss of primary controls could often be flown home "on the trim tabs", or by using trim adjustment as a replacement for the non-working primary controls. Such control is effective, if slower and more limited than primary controls, but it does allow the aircraft to be controlled and directed. In other cases, such as engine failure or damage causing asymmetric drag, trim tabs were invaluable for allowing the pilot to fly the aircraft straight without having to apply a constant force on the stick or rudder to keep the aircraft flying straight.
Trim tabs were also important for aircraft such as bombers, which often underwent rapid changes in center-of-gravity when the bombload was dropped, requiring a hand ready on the trim-adjusting wheel to counteract the tendency of the aircraft to pitch up or down. Undertaking high-speed dives or deploying flaps also generally necessitated pitch trim adjustment, as aircraft of the era had different pitch tendencies at different airspeeds, and flaps could change the center of pressure.
Consumption of fuel could require periodic trim adjustment during a long flight, as it was difficult to ensure that all fuel tanks were equally near the center of gravity. An extreme example was the later P-51 Mustang, which was given a large fuel tank behind the cockpit to allow long-range missions; as fuel from this tank was consumed it was necessary for regular adjustment of the elevator trim.
As a metaphor
Designer Buckminster Fuller is often cited for his use of trim tabs as a metaphor for leadership and personal empowerment. In the February 1972 issue of Playboy, Fuller said:
The official newsletter of the Buckminster Fuller Institute is called Trimtab.
Fuller's metaphor received considerable media attention in January 2019 when actor Jeff Bridges employed it in his Cecil B. DeMille Award acceptance speech at the 76th Golden Globes:
| Technology | Aircraft components | null |
1883900 | https://en.wikipedia.org/wiki/Crab-eating%20fox | Crab-eating fox | The crab-eating fox (Cerdocyon thous), also known as the forest fox, wood fox, bushdog (not to be confused with the bush dog) or maikong, is an extant species of medium-sized canid endemic to the central part of South America since at least the Pleistocene epoch. Like South American foxes, which are in the genus Lycalopex, it is not closely related to true foxes. Cerdocyon comes from the Greek words kerdo (meaning fox) and kyon (dog) referring to the dog- and fox-like characteristics of this animal.
Taxonomy and evolution
The crab-eating fox was originally described as Canis thous by Linnaeus (1766), and first placed in its current genus Cerdocyon by Hamilton-Smith in 1839.
Cerdocyonina is a tribe which appeared around 6.0 million years ago (Mya) in North America as Ferrucyon avius becoming extinct by around 1.4–1.3 Mya. living about . This genus has persisted in South America from an undetermined time, possibly around 3.1 Mya, and continues to the present in the same or a similar form to the crab-eating fox.
As one of the species of the tribe Canini, it is related to the genus Canis. The crab-eating fox's nearest living relative, as theorized at present, is the short-eared dog. This relationship, however, has yet to be supported by mitochondrial investigations. Two subgenera (Atelocynus and Speothos) were long ago included in Cerdocyon.
Cerdocyon thous, C. avius and other species of the genus Cerdocyon underwent radiational evolution on the South American continent. All close relatives of the crab-eating fox (Cerdocyon thous) are extinct. It is the only living representative at present of the genus Cerdocyon.
Description
The crab-eating fox is predominantly greyish-brown, with areas of red on the face and legs, and black-tipped ears and tail. It has short, strong legs and its tail is long and bushy. The head and body length averages , and the average tail length is . It can weigh between .
The coat is short and thick. Coloration varies from grey to brown, to yellowish, to pale, to dark grey. There is a black streak along the back legs, with a black stripe along the spine. On muzzle, ears and paws there is more-reddish fur. The tail, legs and ear tips are black. The ears are wide and round. The torso is somewhat narrow; legs are short but strong. The dense hairy tail stays upright when they are excited. There is significant variation in color between population, from very dark to light grey-yellow.
Genetically, there are 74 diploid chromosomes (36 pairs).
Habitat
The crab-eating fox is a canid that ranges in savannas; woodlands; subtropical forests; prickly, shrubby thickets; and tropical savannas such as the caatinga, plains, and campo, from Colombia and southern Venezuela in the north to Paraguay, Uruguay and northern Argentina at the southernmost reaches of its range. The crab-eating fox has also been sighted in Panama since the 1990s.
Its habitat also includes wooded riverbanks such as riparian forest. In the rainy season, their range moves uphill, whilst in drier times they move to lower ground. Their habitat covers all environments except rainforests, high mountains, and open grassy savannas. In some regions of their range, they are threatened with extirpation.
Behaviour and ecology
The crab-eating fox creates monogamic teams for hunting; groups of several monogamic pairs may form during the reproductive season. Population density estimates vary between one individual per 4 km2 in Venezuela to 0.0003 individuals/km2 in Argentinian wetlands. Territorialism was noticed during the dry season; during rainy seasons, when there is more food, they pay less attention to territory. Hideouts and dens often are found in bushes and in thick grass, and there are typically multiple entrance holes per den. Despite being capable of tunneling, they prefer to take over other animals' burrows. Several characteristic sounds are made by the crab-eating fox such as barking, whirring and howling, which occur often when pairs lose contact with one another.
The crab eating fox is nocturnal, with peaks of activity in the middle of the night and the early morning.
Reproduction
The foxes reach sexual maturity within 9–10 months year. Adult females gives birth to one or two litters per year, depending on the climate and the availability of food. The reproductive period most often begins in November or December, and again in July. The birth of offspring follows after an approximately 56-day gestation, typically in January, February or sometimes March, then again from September to October. If giving birth to one litter, they typically give birth in the early spring. The breeding pair is monogamous and raises the pups together, which are weaned at around three months old and become independent of their parents around 5–8 months old.
Diet
The crab-eating fox searches for crabs on muddy floodplains during the wet season, giving this animal its common name. It is an opportunist and an omnivore, preferring insects or meat from rodents and birds when available. Other foods readily consumed include other crustaceans, tortoises, turtle eggs, bird eggs, insects, lizards, fruit, and carrion. Their diet is varied and has been found to differ by different researchers, suggesting opportunistic feeding and geographical variation. During the wet season, the diet contains more crustaceans, while during the dry season it contains more insects. The crab-eating fox contributes to the control of rodents and harmful insects.
Conservation
The Convention on International Trade in Endangered Species of Wild Fauna and Flora (CITES) lists the fox as not threatened by extinction. The IUCN lists the crab-eating fox as being of "Least Concern". There are no precise estimates of the population size, but it is common within its range and the population is stable.
It is considered a threat to livestock by farmers, which leads to illegal hunting in some countries. The primary threat to the fox is disease from unvaccinated dogs.
Subspecies
The crab-eating fox has five recognized subspecies, differing in sizes and coloring of fur.
C. t. thous, Venezuela, Guyana, Suriname, French Guiana North Brazil.
C. t. azarae, North Brazil.
C. t. entrerianus, Brazil, Bolivia, Uruguay, Paraguay and Argentina.
C. t. aquilus, north Venezuela and Colombia.
C. t. germanus, Bogotá region (Colombia).
| Biology and health sciences | Canines | Animals |
1884115 | https://en.wikipedia.org/wiki/Gall%20wasp | Gall wasp | Gall wasps, also traditionally called gallflies, are hymenopterans of the family Cynipidae in the wasp superfamily Cynipoidea. Their common name comes from the galls they induce on plants for larval development. About 1,300 species of this generally very small creature () are known worldwide, with about 360 species of 36 different genera in Europe and some 800 species in North America.
Features
Like all Apocrita, gall wasps have a distinctive body shape, the so-called wasp waist. The first abdominal tergum (the propodeum) is conjoined with the thorax, while the second abdominal segment forms a sort of shaft, the petiole. The petiole connects with the gaster, which is the functional abdomen in apocritan wasps, starting with the third abdominal segment proper. Together, the petiole and the gaster form the metasoma, while the thorax and the propodeum make up the mesosoma.
The antennae are straight and consist of two or three segments. In many varieties, the backside of the mesosoma appears longitudinally banded. The wings are typically simply structured. The female's egg-depositing ovipositor is often seen protruding from the tip of the metasoma.
Reproduction and development
The reproduction of gall wasps is usually partly parthenogenesis, in which a male is completely unnecessary, and partly two-sex propagation. Most species have alternating generations, with one two-sex generation and one parthenogenic generation annually, whereas some species produce very few males and reproduce only by parthenogenesis, possibly because of infection of the females' gametes by endosymbiotic Wolbachia bacteria. The various generations differentiate both in their appearance and in the form of the plant galls they induce.
The larvae of most gall wasps develop in characteristic plant galls they induce themselves, but many species are instead inquilines of other gall wasps, such as those of the genus Synergus.
The plant galls mostly develop directly after the female insect lays the eggs. The inducement for the gall formation is largely unknown; discussion speculates as to chemical, mechanical, and viral triggers. The hatching larvae nourish themselves with the nutritive tissue of the galls, in which they are otherwise well-protected from external environmental effects. The host plants, and the size and shape of the galls are specific to the majority of gall wasps, with about 70% of the known species parasitizing various types of oak, inducing oak galls. Galls can be found on nearly all parts of such trees, including the leaves, buds, branches, and roots. Other species of gall wasps live in eucalyptus, maple, and many herbs. Species determination is usually much easier through observation of the galls produced rather than of the insect itself.
Parasitism
A gall protects the developing gall wasp for the most vulnerable stage of its life cycle, but many other wasps have found a way to penetrate this defence and parasitise the gall and/or larva(e) within. Some of these inquilines and parasitoids use their long, hardened egg-laying tube (ovipositor) to bore into the gall. These parasitoids may, in turn, be preyed upon by other wasps, hyperparasitoids.
Types
Most species of gall wasps live as gall-formers on oaks. One of the best-known is the common oak gall wasp (Cynips quercusfolii), which induces characteristic, 2-cm in diameter, spherical galls on the undersides of oak leaves.
These turn reddish in the fall and are commonly known as oak apples. Light lentiform galls on the undersides of the same leaves are induced by Neuroterus quercusbaccarum; darker ones with bulging edges are formed by Neuroterus numismalis. Also striking are the galls of Cynips longiventris, which likewise can be found on the undersides of leaves, and are recognizable for their spheroidal shape and irregular red streaks. The oak potato gall wasp (Biorrhiza pallida) has round galls that grow to about 4 cm. These are known colloquially as oak potatoes. The latter type of gall is induced by this type of wasp not on the leaves, but on the roots of the oak. On the buds of young oak twigs, one can often find the hard-shelled galls of Andricus kollari and Andricus quercustozae. Galls do not cause significant harm to oak trees.
Evolution
External phylogeny
The external phylogeny of the Cynipidae is based on Peters et al 2017. The Apocrita is within the "Sawflies" which are shown separately for simplicity here.
Internal phylogeny
The internal phylogeny of gall wasps in the cladogram is based on the molecular phylogenetic analysis of Hearn et al. 2023.
Taxonomy
The Cynipidae contains two subfamilies, one extinct and one extant:
Cynipinae
Hodiernocynipinae†
The Cynipinae consists of nine tribes:
Aulacideini Nieves-Aldrey, Nylander & Ronquist, 2015.
Aylacini Ashmead, 1903.
Ceroptresini Nieves-Aldrey, Nylander & Ronquist, 2015.
Cynipini Billberg, 1820.
Diastrophini Nieves-Aldrey, Nylander & Ronquist, 2015.
Eschatocerini Ashmead, 1903.
Phanacidini Nieves-Aldrey, Nylander & Ronquist, 2015.
Qwaqwaiini Liljeblad, Nieves-Aldrey & Melika, 2011.
Synergini Ashmead, 1896.
In human culture
The galls of several species, especially Mediterranean variants, were once used as tanning agents.
Before his work in human sexuality, Alfred Kinsey was known for his study of gall wasps.
Galls formed on oak trees are one of the main ingredients in iron gall ink.
| Biology and health sciences | Hymenoptera | Animals |
1884746 | https://en.wikipedia.org/wiki/Cortex%20%28botany%29 | Cortex (botany) | In botany, a cortex is an outer layer of a stem or root in a vascular plant, lying below the epidermis but outside of the vascular bundles. The cortex is composed mostly of large thin-walled parenchyma cells of the ground tissue system and shows little to no structural differentiation. The outer cortical cells often acquire irregularly thickened cell walls, and are called collenchyma cells.
Plants
Stems and branches
In the three dimensional structure of herbaceous stems, the epidermis, cortex and vascular cambium form concentric cylinders around the inner cylindrical core of pith. Some of the outer cortical cells may contain chloroplasts, giving them a green color. They can therefore produce simple carbohydrates through photosynthesis.
In woody plants, the cortex is located between the periderm (bark) and the vascular tissue (phloem, in particular). It is responsible for the transportation of materials into the central cylinder of the root through diffusion and may also be used for storage of food in the form of starch.
Roots
In the roots of vascular plants, the cortex occupies a larger portion of the organ's volume than in herbaceous stems. The loosely packed cells of root cortex allow movement of water and oxygen in the intercellular spaces.
One of the main functions of the root cortex is to serve as a storage area for reserve foods. The innermost layer of the cortex in the roots of vascular plants is the endodermis. The endodermis is responsible for storing starch as well as regulating the transport of water, ions and plant hormones.
Lichen
On a lichen, the cortex is also the surface layer or "skin" of the nonfruiting part of the body of some lichens. It is the "skin", or outer layer of tissue, that covers the undifferentiated cells of the . Fruticose lichens have one cortex encircling the branches, even flattened, leaf-like forms. Foliose lichens have different upper and lower cortices. Crustose, placodioid, and squamulose lichens have an upper cortex but no lower cortex, and leprose lichens lack any cortex.
| Biology and health sciences | Plant stem | Biology |
1885646 | https://en.wikipedia.org/wiki/Bennettitales | Bennettitales | Bennettitales (also known as cycadeoids) is an extinct order of seed plants that first appeared in the Permian period and became extinct in most areas toward the end of the Cretaceous. Bennettitales were amongst the most common seed plants of the Mesozoic, and had morphologies including shrub and cycad-like forms. The foliage of bennettitaleans is superficially nearly indistinguishable from that of cycads, but they are distinguished from cycads by their more complex flower-like reproductive organs, at least some of which were likely pollinated by insects.
Although certainly gymnosperms sensu lato (cone-bearing seed plants), the relationships of bennettitaleans to other seed plants is debated. Their general resemblance to cycads is contradicted by numerous more subtle features of their reproductive systems and leaf structure. Some authors have linked bennettitaleans to angiosperms (flowering plants) and gnetophytes (a rare and unusual group of modern gymnosperms), forming a broader group known as Anthophyta. Molecular data contradicts this, with gnetophytes found to be much more genetically similar to conifers. The exact position of Bennettitales remains uncertain.
Description
Bennettitales are divided into two families, Cycadeoidaceae and Williamsoniaceae, which have distinct growth habits. Cycadeoidaceae had stout, cycad-like trunks with bisporangiate (containing both megaspores and microspores) strobili (cones) serving as their reproductive structures. Williamsoniaceae either had bisporangiate or monosporangiate cones, and distinctly slender and branching woody trunks. The Williamsoniaceae grew as woody shrubs with a divaricate branching habit, similar to that of Banksia. It has been suggested that Williamsoniaceae are a paraphyletic (not containing all descendants of a common ancestor) assemblage of all Bennettitales that do not belong to the Cycadeoidaceae.
Foliage
In general, bennettitalean leaves are attached to the stem with a helical (corkscrew) arrangement. Some leaves (most species of Nilssoniopteris, etc.) are narrow, solitary blades with a smooth-edged ("entire") margin. Most leaf morphotypes (Pterophyllum, Ptilophyllum, Zamites, Otozamites, etc.) are pinnate (feather-shaped), with many small leaf segments attached to a central shaft. Others (Anomozamites, a few species of Nilssoniopteris) are incompletely pinnate (sawtooth-shaped) and transitional between these two end members. One unusual leaf form, Eoginkgoites, even approaches a palmate appearance similar to early species of Ginkgo.
The foliage of bennettitaleans resembles that of cycads to such an extent that the foliage of the two groups cannot be reliably distinguished based on gross morphology alone. However, fossil foliage which preserves the cuticle can be assigned to either group with confidence. The stomata of bennettitaleans are described as syndetocheilic. This means that the main paired guard cells develop from the same mother cells as the subsidiary cells which surround them. This contrasts with the haplocheilic stomata of cycads and conifers. In haplocheilic stomata, the ring of subsidiary cells are not derived from the same original structures as the guard cells. This fundamental difference is the main way to differentiate bennettitalean and cycad foliage.
Cones and seeds
Like other gymnosperms, bennettitalean reproductive inflorescences come in the form of cones, which produce pollen and ovules (unfertilized seeds). The cones have a thick central receptacle surrounded by simple, helically-arranged fertile and infertile structures. Tissue at the base of the cone forms layers of scale-like or petal-like bracts to protect the radiating inner structures. Some authors refer to bennettitalean cones as "flowers", though they are not equivalent to true angiosperm flowers. Pollen is often enclosed in paired synangia (pollen sacs). The synangia lie on the adaxial (inner) edge of pollen-bearing leaf-like structures known as microsporophylls. This contrasts with cycads, all of which lack discrete synangia and bear pollen on the abaxial (outer) surface of their microsporophylls.
Many bennettitaleans are bisporangiate, where the pollen and ovules are hosted on the same (bisexual or hermaphrodite) cone. Cavities filled with curved synangia-bearing microsporophylls are encased by thin radiating structures, including thick, infertile interseminal scales and fertile sporophylls with ovules at their tips. The presence of ovules at the tips of sporophylls, rather than the tips of stems, is a major difference between the cones of bennettitaleans and gnetophytes. As the cone is fertilized and matures, the microsporophylls wither away and the ovules transform into seeds.
Most bennettitaleans in the family Williamsoniaceae are instead monosporangiate, with separate pollen and ovule-producing (unisexual) cones on the same plant. The ovule-producing (female) cones (Williamsonia, etc.) are similar to mature bisporangiate cones, with interseminal scales and ovule-tipped sporophylls enclosed by bracts. Pollen-producing (male) cones (Weltrichia, etc.), on the other hand, feature an exposed crown of tapering microsporophylls with adaxial rows of synangia. The microsporophylls may host a single linear row of paired synangia, or instead synangia arranged in a pinnate (feather-shaped) pattern.
Seeds are dicotyledonous (possess two embryonic leaves), with a central embryo surrounded by three layers: the thin megagametophyte, the slightly thicker nucellus, and the protective integument. The upper tip of the seed is tapered and opens through a thin and often extended micropyle. A long, narrow micropyle extending out of the seed is superficially similar to the condition in living gnetophytes. Once the seed is fertilized, the micropyle is sealed by a plug-shaped extension of the nucellus. Unlike living gymnosperms, the tip of the nucellus lacks a pollen chamber (receptacle for stored pollen). The integument is dense and thick, with many layers of differentiated cells. This contrasts with the thin, biseriate (two cell-layer) integument of gnetophytes. Bennettitaleans also lack another gnetophyte-like trait: a sheath of fused bracteoles enveloping the seed. Most integument cells are not unusual in size or shape. However, near the micropyle the innermost layer of integument cells become radially-oriented and elongated, partially closing in on the micropyle. The nucellus and integument are unfused above the chalaza (base of the seed), unlike cycads or gnetophytes, where the layers are fused for much of their height.
Cycadeoidaceans have been suggested to have been self-pollinating, with their stems and cones buried underground, although it has alternatively been proposed that they were pollinated by beetles. The flower-like williamsoniacean male reproductive structure Weltrichia is associated with the female reproductive structure Williamsonia, though it is uncertain whether the parent plants were monoecious (male and female reproductive structures being present on the same plant) or dioecious (where each plant has only one gender of reproductive organ). Weltrichia was likely primarily wind-pollinated, with some species possibly pollinated by beetles.
Several groups of Jurassic and Early Cretaceous insects possessed a long proboscis, and it has been suggested that they fed on nectar produced by bennettitalean reproductive structures, such as the bisexual williamsoniacean reproductive structure Williamsoniella, which had a long, narrow central receptacle which was likely otherwise inaccessible. Early Cretaceous bennettitalean pollen has been found directly associated with a proboscis bearing fly belonging to the extinct family Zhangsolvidae, providing evidence that this family acted as pollinators for the group. The interseminal scales of Bennettitales ovulate cones may have become fleshy at maturity, which could have potentially made then attractive to wild animals that served as seed dispersers.
Taxonomy
History of discovery
The Cycadeoideaceae (originally "Cycadeoideae") were named by English geologist William Buckland in 1828, from fossil trunks found in Jurassic strata on the Isle of Portland, England, which Buckland gave the genus name Cycadeoidea. Buckland provided a description of the family and two species, but failed to give a description of the genus, which has led to Buckland's description of the family being considered invalid by modern taxonomic standards. In publications in 1870, Scottish botanist William Carruthers and English paleobotanist William Crawford Williamson described the first known reproductive organs of the Bennettitales from Jurassic strata of Yorkshire and Jurassic-Cretaceous strata of the Isle of Wight and the Isle of Portland. Caruthers was the first to recognise that Bennettitales had distinct differences from cycads, and established the tribes "Williamsonieae" and "Bennettiteae", with the latter being named after the genus Bennettites named by Caruthers in the same publication, the name being in honour of British botanist John Joseph Bennett. The order Bennettitales was erected by German botanist Adolf Engler in 1892, who recognised the group as separate from the Cycadales.
Relationships to other seed plants
The Anthophyte hypothesis erected by Arber and Parking in 1907 posited that angiosperms arose from Bennettitales, as suggested by the wood-like structures and rudimentary flowers. Based on morphological data, however, Bennettitales were classified as a monophyletic group when paired with Gnetales. a study in 2006 suggested that Bennettitales, Angiosperms, and Gigantopteridales form a clade based on the presence of oleanane. Molecular evidence has consistently contradicted the Anthophyte hypothesis, finding that Angiosperms are the sister group to all living gymnosperms, including Gnetales. Some authors have suggested due to similarities between their seed coats, Bennettitales form a clade with the gymnosperm orders of Gnetales and Erdtmanithecales, dubbed the "BEG group". However, this proposal has been contested by other authors, who contend that these similarities are only superficial and do not indicate a close relationship. A 2017 phylogeny based on molecular signatures of fossilised cuticles found that Bennettitales were more closely related to the Ginkgo+Cycads clade than conifers, and were closely related to Nilssonia and Ptilozamites.
Evolutionary history
The oldest confirmed fossils of bennettitaleans are leaves of Nilssoniopteris shanxiensis, a species from the upper part of the Upper Shihhotse Formation in Shanxi Province, China. This strata is dated to the early Kungurian stage of the early Permian (Cisuralian), around 281 million years ago. Supposed Carboniferous-Permian records of Pterophyllum do not have conclusive bennettitalean affinities or have been reinterpreted as cycad foliage in the form genus Pseudoctenis. True Permian records of benettitalean leaves are rare; outside of the Shihhotse Formation they are only found in the Late Permian (likely Changhsingian)-age Umm Irna Formation in Jordan. This formation is notable for the early occurrence of other Mesozoic-style flora, including the earliest records of corystospermalean foliage (Dicroidium). The order Fredlindiales (containing the genus Fredlindia) from the Late Triassic of Gondwana appears to be closely related to Bennettitales, but differs from it in some aspects of its reproductive organs.
The bennettitalean fossil record reappeared in the Middle Triassic, and williamsoniaceans became globally distributed by the end of the period. The oldest bennettitalean reproductive structures are small Williamsonia "flowers" from the Middle Triassic Esk Formation of Australia. While Williamsoniaceae had a global distribution, Cycadeoidaceae appear to have been primarily confined to the western parts of Laurasia, and are primarily known from the Cretaceous. Bennettitales were widespread and abundant during the Jurassic and Early Cretaceous, however Bennettitales severely declined during the Late Cretaceous, coincident with the rise of flowering plants, being mostly extinct by the end of the period, with the final known remains from the Northern Hemisphere being found in the polar latitude Kakanaut Formation in Chukotka, Russia, dating to the Maastrichtian, assignable to Pterophyllum. A possible late record has been reported from the early Oligocene of eastern Australia and Tasmania, assignable to the genus Ptilophyllum, but no cuticle was preserved, making the referral inconclusive.
Subgroups
Cycadolepis (scales or bracts, unplaced in family)
Haitingeria? (pollen organ?)
Lunzia (pollen organ, unplaced in family)
Leguminanthus? (pollen organ?)
Leuthardtia? (pollen organ?)
Westersheimia (ovulate organ, unplaced in family)
Anthrophyopsis? (leaf)
Family Cycadeoidaceae
Cycadeoidea
Monanthesia
Family Williamsoniaceae
Anomozamites (leaf)
Bennetticarpus (female seed cone)
Bennettistemon (male pollen organ)
Bucklandia (axes)
Eoginkgoites (leaf)
Ischnophyton
Kimuriella (whole plant)
Nilssoniopteris (leaf)
Otozamites (leaf)
Pterophyllum (leaf)
Ptilophyllum (leaf)
Vardekloeftia (female seed cone)
Weltrichia (male pollen organ)
Wielandiella (whole plant)
Williamsonia (female seed cones)
Williamsoniella (bisexual reproductive structure)
Zamites (leaf, in partim)
Haitingeria (pollen organ)
Bennettitales is typically considered the sole order in the class Bennettitopsida Engler (1897) or Cycadeoideopsida Scott (1923). Most paleobotanists prefer the two families as used here, though some authors, such as Anderson & Anderson (2007), classify the order via a larger number of families. Anderson & Anderson also classified the orders Fredlindiales Anderson & Anderson (2003) and Pentoxylales Pilger & Melchior (1954) within Bennettitopsida.
Gallery
| Biology and health sciences | Gymnosperms (except conifers) | Plants |
1886159 | https://en.wikipedia.org/wiki/Microwave%20spectroscopy | Microwave spectroscopy | Microwave spectroscopy is the spectroscopy method that employs microwaves, i.e. electromagnetic radiation at GHz frequencies, for the study of matter.
History
The ammonia molecule NH3 is shaped like a pyramid 0.38 Å in height, with an equilateral triangle of hydrogens forming the base.The nitrogen situated on the axis has two equivalent equilibrium positions above and below the triangle of hydrogens, and this raises the possibility of the nitrogen tunneling up and down, through the plane of the H-atoms. In 1932 Dennison et al. ... analyzed the vibrational energy of this molecule and concluded that the vibrational energy would be split into pairs by the presence of these two equilibrium positions. The next year Wright and Randall observed ... a splitting of 0.67 cm–1 in far infrared lines, corresponding to a frequency of 20 GHz, the value predicted by theory. In 1934 Cleeton and Williams ... constructed a grating echelle spectrometer in order to measure this splitting directly, thereby beginning the field of microwave spectroscopy. They observed a somewhat asymmetric absorption line with a maximum at 24 GHz and a full width at half height of 12 GHz.
In molecular physics
In the field of molecular physics, microwave spectroscopy is commonly used to probe the rotation of molecules.
In condensed matter physics
In the field of condensed matter physics, microwave spectroscopy is used to detect dynamic phenomena of either charges or spins at GHz frequencies (corresponding to nanosecond time scales) and energy scales in the μeV regime. Matching to these energy scales, microwave spectroscopy on solids is often performed as a function of temperature (down to cryogenic regimes of a few K or even lower) and/or magnetic field (with fields up to several T).
Spectroscopy traditionally considers the frequency-dependent response of materials, and in the study of dielectrics microwave spectroscopy often covers a large frequency range. In contrast, for conductive samples as well as for magnetic resonance, experiments at a fixed frequency are common (using a highly sensitive microwave resonator), but frequency-dependent measurements are also possible.
Probing charges in condensed matter physics
For insulating materials (both solid and liquid), probing charge dynamics with microwaves is a part of dielectric spectroscopy.
Amongst the conductive materials, superconductors are a material class that is often studied with microwave spectroscopy, giving information about penetration depth (governed by the superconducting condensate), energy gap (single-particle excitation of Cooper pairs), and quasiparticle dynamics.
Another material class that has been studied using microwave spectroscopy at low temperatures are heavy fermion metals with Drude relaxation rates at GHz frequencies.
Probing spins in condensed matter physics
Microwaves impinging on matter usually interact with charges as well as with spins (via electric and magnetic field components, respectively), with the charge response typically much stronger than the spin response. But in the case of magnetic resonance, spins can be directly probed using microwaves. For paramagnetic materials, this technique is called electron spin resonance (ESR) and for ferromagnetic materials ferromagnetic resonance (FMR).
In the paramagnetic case, such an experiment probes the Zeeman splitting, with a linear relation between the static external magnetic field and the frequency of the probing microwave field. A popular combination, as implemented in commercial X-band ESR spectrometers, is approximately 0.3 T (static field) and 10 GHz (microwave frequency) for a typical material with electron g-factor close to 2.
| Physical sciences | Spectroscopy | Chemistry |
19732090 | https://en.wikipedia.org/wiki/Deforestation%20of%20the%20Amazon%20rainforest | Deforestation of the Amazon rainforest | The Amazon rainforest, spanning an area of 3,000,000 km2 (1,200,000 sq mi), is the world's largest rainforest. It encompasses the largest and most biodiverse tropical rainforest on the planet, representing over half of all rainforests. The Amazon region includes the territories of nine nations, with Brazil containing the majority (60%), followed by Peru (13%), Colombia (10%), and smaller portions in Venezuela, Ecuador, Bolivia, Guyana, Suriname, and French Guiana.
Over one-third of the Amazon rainforest is designated as formally acknowledged indigenous territory, amounting to more than 3,344 territories. Historically, indigenous Amazonian peoples have relied on the forest for various needs such as food, shelter, water, fiber, fuel, and medicines. The forest holds significant cultural and cosmological importance for them. Despite external pressures, deforestation rates are comparatively lower in indigenous territories due to legal land titling initiatives that have reduced deforestation by 75% in Peru.
By the year 2022 around 26% of the forest was considered as deforested or highly degraded. According to the Council on Foreign Relations, 300,000 square miles have been lost.
Cattle ranching and sugar cane growing in the Brazilian Amazon has been identified as the primary cause of deforestation, accounting for about 80% of all deforestation in the region. This makes it the world's largest single driver of deforestation, contributing to approximately 14% of the global annual deforestation. Government tax revenue has subsidized much of the agricultural activity leading to deforestation. By 1995, 70% of previously forested land in the Amazon and 91% of land deforested since 1970 had been converted for cattle ranching. The remaining deforestation primarily results from small-scale subsistence agriculture and mechanized cropland producing crops such as soy and palm.
Satellite data from 2018 revealed a decade-high rate of deforestation in the Amazon, with approximately 7,900 km2 (3,100 sq mi) destroyed between August 2017 and July 2018. The states of Mato Grosso and Pará experienced the highest levels of deforestation during this period. Illegal logging was cited as a cause by the Brazilian environment minister, while critics highlighted the expansion of agriculture as a factor encroaching on the rainforest. Researchers warn that the forest may reach a tipping point where it cannot generate sufficient rainfall to sustain itself. In the first 9 months of 2023 deforestation rate declined by 49.5% due to the policy of Lula's government and international help.
History
In the pre-Columbian era, certain parts of the Amazon rainforest were densely populated and cultivated. However, European colonization in the 16th century, driven by the pursuit of gold and later by the rubber boom, depopulated the region due to diseases and slavery, leading to forest regrowth.
Until the 1970s, access to the largely roadless interior of the forest was challenging, and it remained mostly intact apart from partial clearing along the rivers. Deforestation escalated after the construction of highways penetrating deep into the forest, such as the Trans-Amazonian Highway in 1972.
Challenges arose in parts of the Amazon where poor soil conditions made plantation-based agriculture unprofitable. The crucial turning point in deforestation occurred when colonists began establishing farms within the forest during the 1960s. Their farming practices relied on crop cultivation and the slash-and-burn method. However, due to soil fertility loss and weed invasion, the colonists struggled to effectively manage their fields and crops.
Indigenous areas in the Peruvian Amazon, like the Urarina's Chambira River Basin, experience limited soil productivity, leading to the continual clearing of new lands by indigenous horticulturalists. Cattle raising dominated Amazonian colonization as it required less labor, generated acceptable profits, and involved land under state ownership. While promoted as a reforestation measure, the privatization of land was criticized for potentially encouraging further deforestation and disregarding the rights of Peru's indigenous people, who typically lack formal title to land. The associated law, known as Law 840, faced significant resistance and was eventually repealed as unconstitutional.
Illegal deforestation in the Amazon increased in 2015 after decades of decline, driven primarily by consumer demand for products like palm oil. Brazilian farmers clear land to accommodate the growing demand for crops such as palm oil and soy. Deforestation releases significant amounts of carbon, and if current levels continue, the remaining forests worldwide could disappear within 100 years. The Brazilian government implemented the RED (reducing emissions from deforestation and forest degradation) program to combat deforestation, providing support to various African countries through education programs and financial contributions.
In January 2019, Brazil's president, Jair Bolsonaro, issued an executive order granting the agriculture ministry oversight over certain Amazon lands. This decision has been supported by cattle ranchers and mining companies but criticized for endangering indigenous populations and contributing to Brazil's relative contribution to global climate change.
Reports from the year 2021 indicated a 22% increase in deforestation from the previous year, reaching the highest level since 2006.
Causes of deforestation
The deforestation of the Amazon rainforest is influenced by various factors at local, national, and international levels. The rainforest is sought after for purposes such as cattle ranching, the extraction of valuable hardwoods, land for housing and farming (especially soybeans), the construction of roads (including highways and smaller roads), and the collection of medicinal resources. Deforestation in Brazil is also linked to an economic growth model focused on accumulating factors, primarily land, rather than enhancing overall productivity. It is important to note that illegal logging is a common practice in tree removal during deforestation.
Cattle ranching
According to a 2004 World Bank paper and a 2009 Greenpeace report, cattle ranching in the Brazilian Amazon, supported by the international beef and leather trades, has been identified as responsible for approximately 80% of deforestation in the region. This accounts for about 14% of the world's total annual deforestation, making it the largest driver of deforestation globally. The Food and Agriculture Organization of the United Nations reported in 2006 that 70% of previously forested land in the Amazon, as well as 91% of land deforested since 1970, is now used for livestock pasture.
The 2019 European Union-Mercosur Free Trade Agreement, which establishes one of the world's largest free trade areas, has faced criticism from environmental activists and advocates for indigenous rights. They argue that the trade agreement will contribute to further deforestation of the Amazon rainforest by expanding market access for Brazilian beef.
During Jair Bolsonaro's government, certain environmental laws were weakened, accompanied by reductions in funding and personnel in key government agencies and the dismissal of agency heads and state bodies. The deforestation of the Amazon rainforest accelerated during the COVID-19 pandemic in Brazil. According to Brazil's National Institute for Space Research (INPE), deforestation in the Brazilian Amazon increased by more than 50% in the first three months of 2020 compared to the same period in 2019.
In October 2024, Brazil's environmental protection agency, IBAMA, levied fines totaling 365 million reais (US$64 million) on cattle ranches and meatpacking companies, including JBS SA, the world's largest meat packer, for involvement in illegal deforestation in the Amazon. The fines were imposed on companies accused of raising or purchasing cattle from lands that were deforested without authorization.
Soy bean
Deforestation in the Amazon has occurred as a result of farmers clearing land for mechanized cropland. A study based on NASA satellite data in 2006 revealed that the clearing of land for mechanized cropland had become a significant factor in deforestation in the Brazilian Amazon. This change in land use has had an impact on the region's climate. Researchers discovered that in 2004, a peak year for deforestation, over 20% of the forests in the state of Mato Grosso were converted to cropland. In 2005, when soybean prices decreased by more than 25%, certain areas of Mato Grosso showed a decline in large-scale deforestation events, suggesting that price fluctuations of other crops, beef, and timber could also have a notable influence on future land use in the region.
The cultivation of soybeans, primarily for export and the production of biodiesel and animal feed, has been a significant driver of forest loss in the Amazon. As soybean prices have risen, soy farmers have expanded their activities into forested areas of the Amazon. However, the implementation of a private sector agreement known as the Soy Moratorium has played a crucial role in significantly reducing deforestation associated with soy production in the region. In 2006, several major commodity trading companies, including Cargill, pledged not to purchase soybeans produced in recently deforested areas of the Brazilian Amazon. Prior to the moratorium, 30% of soy field expansion was linked to deforestation, contributing to record-high deforestation rates. After eight years of the moratorium, a study conducted in 2015 found that although the soy production area had expanded by 1.3 million hectares, only about 1% of the new soy expansion had occurred at the expense of forests. In response to the moratorium, farmers opted to plant crops on already cleared land.
The perceived needs of soy farmers have been used to justify certain controversial transportation projects that have been developed in the Amazon. The Belém-Brasília highway (1958) and the Cuiabá-Porto Velho highway (1968) were the only federal highways in the Legal Amazon region that were paved and accessible year-round before the late 1990s. These two highways are considered to be central to the "arc of deforestation," which is presently the primary area of deforestation in the Brazilian Amazon. The Belém-Brasília highway attracted nearly two million settlers in its first twenty years. The success of this highway in opening up the forest was replicated as additional paved roads were constructed, leading to an unstoppable wave of settlement. The completion of these roads was followed by a significant influx of settlers, who also had a substantial impact on the forest.
Logging
Logging in deforestation refers to the practice of cutting down trees for commercial purposes, primarily for the timber industry, which contributes to the overall deforestation of an area. Deforestation is the permanent removal of forests and vegetation cover from an area, often resulting in ecological, social, and economic impacts.
The logging process typically involves the following steps:
Tree selection: Loggers identify and select specific trees for harvesting based on their species, size, and commercial value. Valuable tree species often targeted for logging include mahogany, teak, oak, and other hardwoods.
Access and infrastructure development: Loggers establish infrastructure such as roads and trails within the forest to reach the targeted trees. This infrastructure facilitates the transportation of heavy machinery, logging equipment, and harvested timber.
Clearing vegetation: Prior to logging, loggers often clear the understory vegetation and smaller trees surrounding the target trees to enhance access and maneuverability for machinery.
Tree felling: The selected trees are cut down using chainsaws, harvesters, or other mechanized equipment. The felled trees are then prepared for further processing.
Timber extraction: Once the trees are felled, loggers extract the timber from the forest by removing branches and cutting the tree trunks into logs of appropriate sizes for transport.
Log transportation: Extracted logs are transported from the logging site to processing facilities or storage areas using trucks, barges, or helicopters, depending on the accessibility of the area.
Processing and utilization: At processing facilities, the harvested logs are further processed into lumber, plywood, or other wood products. These products find applications in various industries, such as construction, furniture manufacturing, or paper production.
The impacts of logging on deforestation are significant and wide-ranging.
Loss of biodiversity: Logging often leads to the destruction of forest ecosystems, resulting in the loss of habitat for numerous plant and animal species. Deforestation disrupts the intricate web of biodiversity and can contribute to the extinction or endangerment of various species.
Carbon emissions and climate change: Trees play a crucial role in mitigating climate change by absorbing carbon dioxide through photosynthesis. When trees are logged, the stored carbon is released back into the atmosphere as carbon dioxide, contributing to greenhouse gas emissions and climate change.
Soil erosion and degradation: Forests provide a protective cover for the soil, preventing erosion by wind and water. The removal of trees makes the exposed soil more vulnerable to erosion, leading to the loss of fertile topsoil and the degradation of the land.
Disruption of water cycles: Forests act as natural water catchments, regulating water flow and maintaining water quality. Deforestation can disrupt the water cycle, resulting in reduced water availability, altered rainfall patterns, and an increased risk of droughts or floods.
Indigenous and local community impacts: Many indigenous peoples and local communities depend on forests for their livelihoods, cultural practices, and sustenance. Deforestation and logging can displace these communities, undermine their traditional way of life, and create social conflicts.
Economic considerations: While logging can provide economic benefits in terms of employment and revenue generation, unsustainable logging practices can deplete forest resources and undermine long-term economic sustainability. Overexploitation of forests can lead to the loss of potential future income and economic opportunities.
Efforts to address the impacts of logging on deforestation include implementing sustainable forest management practices, promoting reforestation and afforestation, establishing protected areas, enforcing regulations and policies, and supporting alternative livelihood options for local communities dependent on forests.
A 2013 paper found a correlation between rainforest logging in the Amazon and reduced precipitation in the area, resulting in lower yields per hectare. This suggests that, on a broader scale, there is no economic gain for Brazil through logging, selling trees, and using the cleared land for pastoral purposes.
Oil and Gas Development
Oil and gas projects in the western Amazon are a significant driver of deforestation and associated environmental impacts. These projects contribute to forest loss, water pollution, and the displacement of indigenous peoples. The lack of robust regulatory frameworks exacerbates the vulnerability of these areas to exploitation, creating a multifaceted threat to the Amazon's biodiversity and local communities. According to a September 2016 report by Amazon Watch, the importation of crude oil by the US is linked to about 20,000 sq mi (~50,000 km2) of rainforest destruction in the Amazon and the emission of substantial greenhouse gases. These impacts are mostly focused in the western Amazon countries of Ecuador, Peru, and Colombia. The report also indicates that oil exploration is occurring in an additional ~100,000 sq mi (~250,000 km2) of rainforest.
Roads
95% of the deforestation in the Amazon Rainforest "happens within 3.4 miles of a roadway". Forest clearing always begins near new roads, after expands further. In December 2023 the lower house of the Brazilian Congress approved a bill aiming to pave again the high way BR-319 (Brazil highway), what can threaten the existence of the rainforest. The bill defines the road as “critical infrastructure, indispensable to national security, requiring the guarantee of its trafficability,”
Mining
Mining is a significant contributor to the deforestation of the Amazon rainforest. In the years 2005-2015 it accounted for 9% of deforestation.
Climate change
Climate change seriously increases the likelihood of droughts in the Amazon rainforest. The drought of 2023 for example has been made 30 times more likely. Droughts severely hurt the forest.
Other
During August 2019, a prolonged forest fire occurred in the Amazon, contributing significantly to deforestation during that summer. Approximately 519 sq mi (1,340 km2) of the Amazon forest was lost.
It is worth noting that certain instances of deforestation in the Amazon have been attributed to farmers clearing land for small-scale subsistence agriculture.
Loss rates
During the early 2000s, deforestation in the Amazon rainforest showed an increasing trend, with an annual rate of 27,423 km2 (10,588 sq mi) of forest loss recorded in 2004. Subsequently, the annual rate of forest loss generally slowed between 2004 and 2012, although there were spikes in deforestation rates in 2008, 2013, and 2015.
However, recent data suggests that the loss of forest cover is once again accelerating. Between August 2017 and July 2018, approximately 7,900 km2 (3,100 sq mi) of forest were deforested in Brazil, representing a 13.7% increase compared to the previous year and the largest area cleared since 2008. Deforestation in the Brazilian Amazon rainforest experienced a significant surge in June 2019, rising more than 88% compared to the same month in 2018. and more than doubling in January 2020 compared to January 2019.
In August 2019, a substantial number of forest fires, totaling 30,901 individual fires, were reported, marking a threefold increase compared to the previous year. However, the number of fires decreased by one-third in September, and by October 7, it had dropped to approximately 10,000. It is important to note that deforestation is considered to have more severe consequences than burning. The National Institute for Space Research (INPE) in Brazil estimated that at least 7,747 km2 (2,991 sq mi) of the Brazilian Amazon rainforest were cleared during the first half of 2019. INPE subsequently reported that deforestation in the Brazilian Amazon reached a 12-year high between August 2019 and July 2020.
Deforestation figures in Brazil are annually provided by the Instituto Nacional de Pesquisas Espaciais (INPE), based on satellite images captured during the dry season in the Amazon by the Landsat satellite. It's important to note that these estimates may focus solely on the loss of the Amazon rainforest and may not include the loss of natural fields or savannah within the Amazon biome.
Estimated loss by year
†Value calculated from estimated forest loss, not directly known.
Impacts
Deforestation and biodiversity loss in the Amazon rainforest have resulted in significant risks of irreversible changes. Modeling studies have suggested that deforestation may be approaching a critical "tipping point" where large-scale "savannization" or desertification could occur, leading to catastrophic consequences for the global climate. This tipping point could trigger a self-perpetuating collapse of biodiversity and ecosystems in the region. Failing to prevent this tipping point could have severe impacts on the economy, natural capital, and ecosystem services. A study published in Nature Climate Change in 2022 provided empirical evidence that more than three-quarters of the Amazon rainforest has experienced a decline in resilience since the early 2000s, posing risks of dieback that would impact biodiversity, carbon storage, and climate change.
To maintain a high level of biodiversity, research suggests that a threshold of 40% forest cover in the Amazon should be maintained.
Impact on global warming
Deforestation, along with other forms of ecosystem destruction such as peatbog degradation, can have multiple effects. It can reduce the carbon sink capacity of the land and contribute to increased emissions through factors like wildfires, land-use change, and reduced ecosystem health. These impacts can disrupt the normal carbon-absorbing processes of ecosystems, leading to stress and imbalance.
Historically, the Amazon Basin has played a significant role as a carbon sink, absorbing approximately 25% of the carbon captured by terrestrial land.
However, a scientific review article published in 2021 indicates that current evidence suggests the Amazon basin is currently emitting more greenhouse gases than it absorbs overall. This shift is attributed to climate change impacts and human activities in the region, particularly wildfires, current land-use practices, and deforestation. These factors contribute to the release of forcing agents that are likely to result in a net warming effect. Warming temperatures and changing weather patterns also lead to physiological responses in the forest, further hindering the absorption of .
According to a research led by Elena Shevliakova and Stephen Pacala complete deforestation of the Amazon will cause a global temperature rise of 0.25 degrees. In 2023, despite deforestation, the forest still held more than 150 billion metric tons of carbon.
Impacts on water supply
The deforestation of the Amazon rainforest has had a significant impact on Brazil's freshwater supply, particularly affecting the agricultural industry, which has been involved in clearing the forests. In 2005, certain regions of the Amazon basin experienced the most severe drought in over a century. This can be attributed to two key factors:
1. The rainforest plays a crucial role in contributing to rainfall across Brazil, even in distant areas. Deforestation has exacerbated the effects of droughts in 2005, 2010, and 2015–2016.
2. The rainforest contributes to rainfall and facilitates water storage, which in turn provides freshwater to the rivers that supply Brazil and other countries with water.
Impact on local temperature
In 2019, a group of scientists conducted research indicating that under a "business as usual" scenario, the deforestation of the Amazon rainforest will lead to a temperature increase of 1.45 degrees in Brazil. They stated that this temperature rise could have various consequences, including increased human mortality rates and electricity demands, reduced agricultural yields and water resources, and the potential collapse of biodiversity, especially in tropical regions. Additionally, local warming may cause shifts in species distributions, including those involved in the transmission of infectious disease. The authors of the paper assert that deforestation is already contributing to the observed temperature rise.
According to another research, complete Amazon deforestation will render the region itself (including 7 million square kilometers, 9 states in Brazil and 8 other countries) more or less uninhabitable as the temperature will rise by more than 4.5 degrees and rainfall will be reduced by a quart.
Impact on indigenous people
More than one-third of the Amazon forest is designated as Indigenous territory, encompassing over 4,466 formally recognized territories. Until 2015, approximately 8% of deforestation in the Amazon occurred within forests inhabited by indigenous peoples, while 88% occurred in areas outside of indigenous territories and protected areas, despite these areas comprising less than 50% of the total Amazon region. Indigenous communities have historically relied on the forest for sustenance, shelter, water, materials, fuel, and medicinal resources. The forest holds significant cultural and spiritual importance for them. Consequently, deforestation rates tend to be lower within Indigenous Territories, even though pressures to clear land for other purposes persist.
During the deforestation of the Amazon, native tribes have often faced mistreatment and abuse. Encroachments by loggers onto indigenous lands have led to conflicts resulting in fatalities. Some uncontacted indigenous groups have emerged from the forests and interacted with mainstream society due to threats from outsiders. When uncontacted tribes come into contact with outsiders, they are vulnerable to diseases against which they have little immunity. As a result, entire tribes can be severely impacted by epidemics, leading to significant population declines within a few years.
A long-standing struggle has taken place over the control of indigenous territories in the Amazon, primarily involving the Brazilian government. The demand for these lands has stemmed, in part, from the aim of enhancing Brazil's economic standing. Various individuals, including ranchers and land speculators from the southeast, have sought to claim these lands for personal financial gain. In early 2019, Brazil's newly elected president, Jair Bolsonaro, issued an executive order empowering the agriculture ministry to regulate the land occupied by indigenous tribes in the Amazon.
In the past, mining operations were permitted within the territory of an isolated indigenous group called the Yanomami. The conditions endured by these indigenous peoples resulted in many health issues, including tuberculosis. If their lands are utilized for further development, numerous tribal communities will be forcibly displaced, potentially leading to the loss of lives. In addition to the mistreatment of indigenous peoples, the exploitation of the forest itself will result in the depletion of vital resources necessary for their daily lives.
Research conducted in the Peruvian Amazon demonstrates that legal titling of indigenous lands significantly reduces deforestation rates. Land titling initiatives led to a 75% reduction in deforestation over two years. Such policies provide indigenous communities with the legal authority to protect their land from encroachment and exploitation, creating a powerful tool for combating deforestation in the Amazon.
Between 2013 and 2021, deforestation within indigenous territories in the Brazilian Amazon increased by 129%, primarily driven by illegal mining activities, activities which threaten biodiversity and undermine the cultural and environmental integrity of these lands.
The Yanomami people in Brazil have faced severe challenges due to illegal gold mining in their territories. Mining activities have led to deforestation, water contamination, and a surge in malaria cases among the Yanomami. President Lula's administration has taken steps to address these issues, including efforts to reclaim and protect Yanomami lands. In 2012, the Yanomami organization HORONAMI expressed their concerns, stating, “Illegal miners persist in destroying our lands, and the government must act urgently to stop these abuses and investigate the harm being done in the Upper Ocamo region”. In 2020, a report from Mongabay quoted the Yanomami leaders expressing frustration with Brazilian authorities, stating, “We feel utterly abandoned; the government turns a blind eye to the illegal activities that poison our rivers and bring disease to our people”.
Dario Kopenawa, vice president of the Hutukara Yanomami Association, has emphasized the importance of government intervention, stating, “The Brazilian government must fulfil its protective role, where every Brazilian citizen, not just the Yanomami, feels protected. It is not a favour, but a constitutional obligation. It is necessary to curb the mining projects on indigenous lands because they are illegal under Brazilian law.”
Recent government raids targeting illegal gold mining in the Yanomami Indigenous Territory have revealed the extreme extent of deforestation caused by these activities. Using machine learning algorithms and satellite imagery, researchers estimate that over 2,000 hectares of forest have been deforested due to gold mining since 2019, with 67% of this deforestation (approximately 1,350 hectares) occurring in 2022 alone. The deforestation is widespread, affecting areas along the Uraricoera, Parima, and Mucajai Rivers.
Efforts to stop and reverse deforestation
Norwegian Prime Minister Jens Stoltenberg announced on September 16, 2008, that the Norwegian government would contribute a donation of US$1 billion to the newly established Amazon Fund. The funds from this initiative would be allocated to projects aimed at mitigating the deforestation of the Amazon rainforest.
In September 2015, Brazilian President Dilma Rousseff addressed the United Nations, reporting that Brazil had effectively reduced the deforestation rate in the Amazon by 82%. She also outlined Brazil's goals for the next 15 years, which included eliminating illegal deforestation, restoring and reforesting an area of 120,000 km2 (46,000 sq mi), and rehabilitating 150,000 km2 (58,000 sq mi) of degraded pastures.
In August 2017, Brazilian President Michel Temer revoked the protected status of an Amazonian nature reserve, which spanned an area equivalent to Denmark in the northern states of Pará and Amapá.
In April 2019, an Ecuadorian court issued an order to cease oil exploration activities in an area of 1,800 square kilometers (690 sq mi) within the Amazon rainforest.
In May 2019, eight former environment ministers in Brazil expressed concerns about escalating deforestation in the Amazon during Jair Bolsonaro's first year as president. Carlos Nobre, an expert on the Amazon and climate change, warned in September 2019 that if deforestation rates continued at their current pace, the Amazon forest could reach a tipping point within 20 to 30 years, potentially resulting in large portions of the forest transforming into a dry savanna, particularly in the southern and northern regions.
Bolsonaro has rebuffed European politicians' attempts to intervene in the matter of Amazon rainforest deforestation, citing it as Brazil's internal affairs. He has advocated for the opening of more areas, including those in the Amazon, for mining activities and mentioned discussions with US President Donald Trump about a joint development program for the Brazilian Amazon region.
Brazilian Economy Minister Paulo Guedes has expressed the belief that other countries should compensate Brazil for the oxygen produced within its borders but used elsewhere.
In late August 2019, following an international outcry and warnings from experts about the escalating fires, the Brazilian government, led by Jair Bolsonaro, implemented measures to combat the fires. These measures included a 60-day ban on forest clearance using fires, deploying 44,000 soldiers to fight the fires, receiving four planes from Chile for firefighting purposes, accepting a $12 million aid package from the UK government, and softening his stance on aid from the G7. Bolsonaro also called for a Latin American conference to address the preservation of the Amazon.
On November 2, 2021, during the COP26 climate summit, over 100 countries, representing approximately 85% of the world's forests, reached a significant agreement to end deforestation by 2030. This agreement, an improvement on the 2014 New York Declaration on Forests, which initially aimed to reduce deforestation by 50% by 2020 and end it by 2030, now includes Brazil as a signatory. It is worth noting that deforestation increased during the 2014–2020 period despite the previous agreement.
In August 2023, Brazilian President Luiz Inácio Lula da Silva hosts a summit in Belem with eight South American countries to coordinate policies for the Amazon basin and develop a roadmap to save the world's largest rainforest, also serving as a preparatory event for the COP30 UN climate talks in 2025.
In the first 8 months of 2023 deforestation rate in the Brazilian Amazon declined by 48%, that prevented the release of 196 million tons CO2 to the atmosphere. Financing from the Amazon Fund and cooperation between the Amazonian nations played a significant role in it. In the first 9 months of 2023 deforestation rate declined by 49.5% despite the worst drought in the last 40 years. Wildfires in September 2023 declined by 36% in comparison to September 2022. Switzerland and United States gave 8.4 million dollars to the Amazon fund for preventing deforestation.
According to Amazon Conservation's MAAP forest monitoring program, the deforestation rate in the Amazon rainforest as a whole from the 1 of January to the 8 of November 2023, in comparison to the same period in 2022, declined by 55.8%. This trend gives hope to Amazon. Deforestation reduction in Brazil (59%) which is the main cause of the trend is probably due to Lula's environmental policy. In Columbia, the rate of deforestation fell by 66.5% probably due to the policies of Gustavo Petro and a change in the policies of former guerrilla fighters that control part of the forest. It is not clear, still, what caused the decline in Bolivia (60%) and Peru (37%). Bolivia has a relatively high forest loss rate due to wildfires, but those are not occurring in the Amazon.
In September 2024, Sawré Muybu, an indigenous land, belonging to the Munduruku people got an official recognition, which is considered as a significant step in fighting deforestation. However, 44 more territories still wait for recognition.
Cost of rainforest conservation
According to the Woods Hole Research Institute (WHRC) in 2008, it was estimated that halting deforestation in the Brazilian rainforest would require an annual investment of US$100–600 million. A more recent study in 2022 suggested that the conservation of approximately 80% of the Brazilian rainforest remains achievable, with an estimated annual cost of US$1.7–2.8 billion to conserve an area of 3.5 million km2. By preventing deforestation, it would be possible to avoid carbon emissions at a cost of US$1.33 per ton of , which is lower compared to the cost of reducing emissions through renewable fuel subsidies (US$100 per ton) or weatherization assistance programs such as building insulation (US$350/ton).
Future of the Amazon rainforest
Based on the deforestation rates observed in 2005, projections indicated that the Amazon rainforest would experience a 40% reduction within two decades. While the rate of deforestation has slowed since the early 2000s, the forest continues to shrink annually, and satellite data analysis reveals a significant increase in deforestation since 2018.
| Physical sciences | Climate change | Earth science |
19740545 | https://en.wikipedia.org/wiki/Traffic%20collision | Traffic collision | A traffic collision, also known as a motor vehicle collision, or car crash, occurs when a vehicle collides with another vehicle, pedestrian, animal, road debris, or other moving or stationary obstruction, such as a tree, pole or building. Traffic collisions often result in injury, disability, death, and property damage as well as financial costs to both society and the individuals involved. Road transport is statistically the most dangerous situation people deal with on a daily basis, but casualty figures from such incidents attract less media attention than other, less frequent types of tragedy. The commonly used term car accident is increasingly falling out of favor with many government departments and organizations, with the Associated Press style guide recommending caution before using the term and the National Union of Journalists advising against it in their Road Collision Reporting Guidelines. Some collisions are intentional vehicle-ramming attacks, staged crashes, vehicular homicide or vehicular suicide.
Several factors contribute to the risk of collisions, including vehicle design, speed of operation, road design, weather, road environment, driving skills, impairment due to alcohol or drugs, and behavior, notably aggressive driving, distracted driving, speeding and street racing.
In 2013, 54 million people worldwide sustained injuries from traffic collisions. This resulted in 1.4 million deaths in 2013, up from 1.1 million deaths in 1990. About 68,000 of these occurred with children less than five years old. Almost all high-income countries have decreasing death rates, while the majority of low-income countries have increasing death rates due to traffic collisions. Middle-income countries have the highest rate with 20 deaths per 100,000 inhabitants, accounting for 80% of all road fatalities with 52% of all vehicles. While the death rate in Africa is the highest (24.1 per 100,000 inhabitants), the lowest rate is to be found in Europe (10.3 per 100,000 inhabitants).
Terminology
Traffic collisions can be classified by general types. Types of collision include head-on, road departure, rear-end, side collisions, and rollovers.
Many different terms are commonly used to describe vehicle collisions. The World Health Organization uses the term road traffic injury, while the U.S. Census Bureau uses the term motor vehicle accidents (MVA), and Transport Canada uses the term "motor vehicle traffic collision" (MVTC). Other common terms include auto accident, car accident, car crash, car smash, car wreck, motor vehicle collision (MVC), personal injury collision (PIC), road accident, road traffic accident (RTA), road traffic collision (RTC), and road traffic incident (RTI) as well as more unofficial terms including smash-up, pile-up, and fender bender
Criticism of "accident" terminology
Many organizations, companies and government agencies have begun to avoid the term accident, instead preferring terms such as collision, crash or incident. This is because the term accident may imply that there is no one to blame or that the collision was unavoidable, whereas most traffic collisions are the result of driving under the influence, excessive speed, distractions such as mobile phones, other risky behavior, poor road design, or other preventable factors.
In 1997, George L. Reagle, the Associate Administrator for Motor Carriers of the Federal Motor Carrier Safety Administration wrote a letter stating that "A crash is not an accident", emphasizing that the Department's Research and Special Programs Administration, the Federal Highway Administration, and the National Highway Traffic Safety Administration had all declared that "accident" should be avoided in their published writings and media communications. In 2016, the Associated Press updated its style guide to recommend that journalists use "crash, collision, or other terms" rather than "accident" unless culpability is proven. The AP also recommends avoiding "accident" when negligence is proven or claimed because the term "can be read as exonerating the person responsible." In 2021, the American Automobile Association (AAA) passed a resolution to replace "car accident" with "car crash" in their vocabulary. In 2022, the traffic management company INRIX announced that "accident" would be removed from their lexicon. In 2023, the National Union of Journalists in the UK published the Road Collision Reporting Guidelines which includes a recommendation that journalists should "Avoid use of the word ‘accident’ until the facts of a collision are known."
The Maryland Department of Transportation's Highway Safety Office emphasizes that "crashes are no accident", saying that "Using the word accident suggests that an incident was unavoidable, but many roadway crashes can be attributed to human error." The Michigan Department of Transportation states that "accident" should be dropped in favor of "crash", saying that "Traffic crashes are fixable problems, caused by inattentive drivers and driver behavior. They are not accidents." In line with their Vision Zero commitments, the Portland Bureau of Transportation recommends using "crash" rather than "accident".
On the contrary, some have criticized the use of terminology other than accident for holding back safety improvements, based on the idea that such terms perpetuate a culture of blame that may discourage the involved parties from fully disclosing the facts, and thus frustrate attempts to address the real root causes.
Intent
Some traffic collisions are caused intentionally by a driver. For example, a collision may be caused by a driver who intends to commit vehicular suicide. Collisions may also be intentionally caused by people who hope to make an insurance claim against the other driver or may be staged for such purposes as insurance fraud. Motor vehicles may also be involved in collisions as part of a deliberate effort to hurt other people, such as in a vehicle-ramming attack or vehicular homicide.
Health effects
Physical
A number of physical injuries can commonly result from the blunt force trauma caused by a collision, ranging from bruising and contusions to catastrophic physical injury (e.g., paralysis), traumatic or non-traumatic cardiac arrest and death. The CDC estimates that roughly 100 people die in motor vehicle crashes each day in the United States.
Psychological
Following collisions, long-lasting psychological trauma may occur. These issues may make those who have been in a crash afraid to drive again. In some cases, psychological trauma may affect individuals' lives, causing difficulty going to work, attending school, or performing family responsibilities.
Causes
Road incidents are caused by a large number of human factors such as failing to act according to weather conditions, road design, signage, speed limits, lighting conditions, pavement markings, and roadway obstacles. A 1985 study by K. Rumar, using British and American crash reports as data, suggested 57% of crashes were due solely to driver factors, 27% to the combined roadway and driver factors, 6% to the combined vehicle and driver factors, 3% solely to roadway factors, 3% to combined roadway, driver, and vehicle factors, 2% solely to vehicle factors, and 1% to combined roadway and vehicle factors. Reducing the severity of injury in crashes is more important than reducing incidence and ranking incidence by broad categories of causes is misleading regarding severe injury reduction. Vehicle and road modifications are generally more effective than behavioral change efforts with the exception of certain laws such as required use of seat belts, motorcycle helmets, and graduated licensing of teenagers.
Human factors
Human factors in vehicle collisions include anything related to drivers and other road users that may contribute to a collision. Examples include driver behavior, visual and auditory acuity, decision-making ability, and reaction speed.
A 1985 report based on British and American crash data found driver error, intoxication, and other human factors contribute wholly or partly to about 93% of crashes. A 2019 report from the U.S. National Highway Traffic Safety Administration found that leading contributing factors for fatal crashes included driving too fast for conditions or in excess of the speed limit, operating under the influence, failure to yield right of way, failure to keep within the proper lane, operating a vehicle in a careless manner, and distracted driving.
Drivers distracted by mobile devices had nearly four times greater risk of crashing their cars than those who were not. Research from the Virginia Tech Transportation Institute has found that drivers who are texting while driving are 23 times more likely to be involved in a crash as non-texting drivers. Dialing a phone is the most dangerous distraction, increasing a drivers' chance of crashing by 12 times, followed by reading or writing, which increased the risk by ten times.
An RAC survey of British drivers found 78% of drivers thought they were highly skilled at driving, and most thought they were better than other drivers, a result suggesting overconfidence in their abilities. Nearly all drivers who had been in a crash did not believe themselves to be at fault. One survey of drivers reported that they thought the key elements of good driving were:
controlling a car including a good awareness of the car's size and capabilities
reading and reacting to road conditions, weather, road signs, and the environment
alertness, reading and anticipating the behavior of other drivers.
Although proficiency in these skills is taught and tested as part of the driving exam, a "good" driver can still be at a high risk of crashing because:
An Axa survey concluded Irish drivers are very safety-conscious relative to other European drivers. This does not translate to significantly lower crash rates in Ireland.
Accompanying changes to road designs have been wide-scale adoptions of rules of the road alongside law enforcement policies that included drink-driving laws, setting of speed limits, and speed enforcement systems such as speed cameras. Some countries' driving tests have been expanded to test a new driver's behavior during emergencies, and their hazard perception.
There are demographic differences in crash rates. For example, although young people tend to have good reaction times, disproportionately more young male drivers feature in collisions, with researchers observing that many exhibit behaviors and attitudes to risk that can place them in more hazardous situations than other road users. This is reflected by actuaries when they set insurance rates for different age groups, partly based on their age, sex, and choice of vehicle. Older drivers with slower reactions might be expected to be involved in more collisions, but this has not been the case as they tend to drive less and, apparently, more cautiously. Attempts to impose traffic policies can be complicated by local circumstances and driver behavior. In 1969 Leeming warned that there is a balance to be struck when "improving" the safety of a road.
Conversely, a location that does not look dangerous may have a high crash frequency. This is, in part, because if drivers perceive a location as hazardous, they take more care. Collisions may be more likely to happen when hazardous road or traffic conditions are not obvious at a glance, or where the conditions are too complicated for the limited human machine to perceive and react in the time and distance available. High incidence of crashes is not indicative of high injury risk. Crashes are common in areas of high vehicle congestion, but fatal crashes occur disproportionately on rural roads at night when traffic is relatively light.
This phenomenon has been observed in risk compensation research, where the predicted reductions in collision rates have not occurred after legislative or technical changes. One study observed that the introduction of improved brakes resulted in more aggressive driving, and another argued that compulsory seat belt laws have not been accompanied by a clearly attributed fall in overall fatalities. Most claims of risk compensation offsetting the effects of vehicle regulation and belt use laws have been discredited by research using more refined data.
In the 1990s, Hans Monderman's studies of driver behavior led him to the realization that signs and regulations had an adverse effect on a driver's ability to interact safely with other road users. Monderman developed shared space principles, rooted in the principles of the woonerven of the 1970s. He concluded that the removal of highway clutter, while allowing drivers and other road users to mingle with equal priority, could help drivers recognize environmental clues. They relied on their cognitive skills alone, reducing traffic speeds radically and resulting in lower levels of road casualties and lower levels of congestion.
Some crashes are intended; staged crashes, for example, involve at least one party who hopes to crash a vehicle in order to submit lucrative claims to an insurance company. In the United States during the 1990s, criminals recruited Latin American immigrants to deliberately crash cars, usually by cutting in front of another car and slamming on the brakes. It was an illegal and risky job, and they were typically paid only $100. Jose Luis Lopez Perez, a staged crash driver, died after one such maneuver, leading to an investigation that uncovered the increasing frequency of this type of crash.
Motor vehicle speed
The U.S. Department of Transportation's Federal Highway Administration reviewed research on traffic speed in 1998. The summary says that:
The evidence shows the risk of having a crash is increased both for vehicles traveling slower than the average speed and for those traveling above the average speed.
The risk of being injured increases exponentially with speeds much faster than the median speed.
The severity/lethality of a crash depends on the vehicle speed change at impact.
There is limited evidence suggesting lower speed limits result in lower speeds on a system-wide basis.
Most crashes related to speed involve speed too fast for the conditions.
More research is needed to determine the effectiveness of traffic calming.
In the U.S. in 2018, 9,378 people were killed in motor vehicle crashes involving at least one speeding driver, which accounted for 26% of all traffic-related deaths for the year.
In Michigan in 2019, excessive speed was a factor in 18.8% of the fatalities that resulted from fatal motor vehicle crashes and in 15.6% of the suspected serious injuries resulting from crashes.
The Road and Traffic Authority (RTA) of the Australian state of New South Wales (NSW) asserts speeding (traveling too fast for the prevailing conditions or above the posted speed limit) is a factor in about 40 percent of road deaths. The RTA also says speeding increases the risk of a crash and its severity. On another web page, the RTA qualifies its claims by referring to one specific piece of research from 1997, and writes "Research has shown that the risk of a crash causing death or injury increases rapidly, even with small increases above an appropriately set speed limit."
The contributory factor report in the official British road casualty statistics shows for 2006, that "exceeding the speed limit" was a contributory factor in 5% of all casualty crashes (14% of all fatal crashes), and "traveling too fast for conditions" was a contributory factor in 11% of all casualty crashes (18% of all fatal crashes).
In France, in 2018, the speed limit was reduced from 90 km/h to 80 km/h on a large part of the local outside built-up area road network in the sole aim of reducing the number of road fatalities.
Assured clear distance ahead
A common cause of collisions is driving faster than one can stop within their field of vision. Such practice is illegal and is particularly responsible for an increase in fatalities at night – when it occurs most.
Driver impairment
Driver impairment describes factors that prevent the driver from driving at their normal level of skill. Common impairments include:
Alcohol
According to the Government of Canada, coroner reports from 2008 suggested almost 40% of fatally injured drivers consumed some quantity of alcohol before the collision.
Physical impairment
Poor eyesight and/or physical impairment, with many jurisdictions setting simple sight tests and/or requiring appropriate vehicle modifications before being allowed to drive.
Youth
Insurance statistics demonstrate a notably higher incidence of collisions and fatalities among drivers aged in their teens or early twenties, with insurance rates reflecting this data. These drivers have the highest incidence of both collisions and fatalities among all driver age groups, a fact that was observed well before the advent of mobile phones.Females in this age group exhibit somewhat lower collision and fatality rates than males but still register well above the median for drivers of all ages. Also within this group, the highest collision incidence rate occurs within the first year of licensed driving. For this reason, many US states have enacted a zero-tolerance policy wherein receiving a moving violation within the first six months to one year of obtaining a license results in automatic license suspension. South Dakota is the only state that allows fourteen-year-olds to obtain drivers' licenses.
Old age
Old age, with some jurisdictions requiring driver retesting for reaction speed and eyesight after a certain age.
Sleep deprivation
Various factors such as fatigue or sleep deprivation might increase the risk, or the number of hours of driving might increase the risk of an incident. 41% of drivers self-report having fallen asleep at the wheel. It is estimated that 15% of fatal crashes involve drowsiness (10% of daytime crashes, and 24% of nighttime crashes). Work factors can increase the risk of drowsy driving such as long or irregular hours or driving at night.
Drug use
Including some prescription drugs, over-the-counter drugs (notably antihistamines, opioids and muscarinic antagonists), and illegal drugs.
Distraction
Research suggests that the driver's attention is affected by distracting sounds such as conversations and operating a mobile phone while driving. Many jurisdictions now restrict or outlaw the use of some types of phones in the car. Recent research conducted by British scientists suggests that music can also have an effect; classical music is considered to be calming, yet too much could relax the driver to a condition of distraction. On the other hand, hard rock may encourage the driver to step on the acceleration pedal, thus creating a potentially dangerous situation on the road.Cell phone use is an increasingly significant problem on the roads and the U.S. National Safety Council compiled more than 30 studies postulating that hands-free is not a safer option because the brain remains distracted by the conversation and cannot focus solely on the task of driving.
Combinations
Several conditions can combine to create a more dangerous situation, for example, low doses of alcohol and cannabis have a more severe effect on driving performance than either in isolation. Taking recommended doses of several drugs together, which individually do not cause impairment, may cause drowsiness. This could be more pronounced in an elderly person whose renal function is less efficient than a younger person's.
Road design
A 1985 US study showed that about 34% of serious crashes had contributing factors related to the roadway or its environment. Most of these crashes also involved a human factor. The road or environmental factor was either noted as making a significant contribution to the circumstances of the crash or did not allow room to recover. In these circumstances, it is frequently the driver who is blamed rather than the road; those reporting the collisions have a tendency to overlook the human factors involved, such as the subtleties of design and maintenance that a driver could fail to observe or inadequately compensate for.
Research has shown that careful design and maintenance, with well-designed intersections, road surfaces, visibility and traffic control devices, can result in significant improvements in collision rates.
Individual roads also have widely differing performance in the event of an impact. In Europe, there are now EuroRAP tests that indicate how "self-explaining" and forgiving a particular road and its roadside would be in the event of a major incident.
In the UK, research has shown that investment in a safe road infrastructure program could yield a reduction in road deaths, saving as much as £6 billion per year. A consortium of 13 major road safety stakeholders has formed the Campaign for Safe Road Design, which is calling on the UK Government to make safe road design a national transport priority.
Vehicle design and maintenance
Seat belts
Research has shown that, across all collision types, it is less likely that seat belts were worn in collisions involving death or serious injury, rather than light injury; wearing a seat belt reduces the risk of death by about 45 percent. Seat belt use is controversial, with notable critics such as Professor John Adams suggesting that their use may lead to a net increase in road casualties due to a phenomenon known as risk compensation. Observation of driver behaviors before and after seat belt laws does not support the risk compensation hypothesis.
Several driving behaviors were observed on the road before and after the belt use law was enforced in Newfoundland, and in Nova Scotia during the same period without a law. Belt use increased from 16 percent to 77 percent in Newfoundland and remained virtually unchanged in Nova Scotia. Four driver behaviors (speed, stopping at intersections when the control light was amber, turning left in front of oncoming traffic, and gaps in following distance) were measured at various sites before and after the law. Changes in these behaviors in Newfoundland were similar to those in Nova Scotia, except that drivers in Newfoundland drove slower on expressways after the law, contrary to the risk compensation theory.
Maintenance
A well-designed and well-maintained vehicle, with good brakes, tires and well-adjusted suspension will be more controllable in an emergency and thus be better equipped to avoid collisions. Some mandatory vehicle inspection schemes include tests for some aspects of roadworthiness, such as the UK's MOT test or German TÜV conformance inspection.
The design of vehicles has also evolved to improve protection after collision, both for vehicle occupants and for those outside of the vehicle. Much of this work was led by automotive industry competition and technological innovation, leading to measures such as Saab's safety cage and reinforced roof pillars of 1946, Ford's 1956 Lifeguard safety package, and Saab and Volvo's introduction of standard fit seatbelts in 1959. Other initiatives were accelerated as a reaction to consumer pressure, after publications such as Ralph Nader's 1965 book Unsafe at Any Speed accused motor manufacturers of indifference to safety.
In the early 1970s, British Leyland started an intensive program of vehicle safety research, producing a number of prototype experimental safety vehicles demonstrating various innovations for occupant and pedestrian protection such as airbags, anti-lock brakes, impact-absorbing side-panels, front and rear head restraints, run-flat tires, smooth and deformable front-ends, impact-absorbing bumpers, and retractable headlamps. The design has also been influenced by government legislation, such as the Euro NCAP impact test.
Common features designed to improve safety include thicker pillars, safety glass, interiors with no sharp edges, stronger bodies, other active or passive safety features, and smooth exteriors to reduce the consequences of an impact on pedestrians.
The UK Department for Transport publish road casualty statistics for each type of collision and vehicle through its Road Casualties Great Britain report.
These statistics show a ten-to-one ratio of in-vehicle fatalities between types of cars. In most cars, occupants have a 2–8% chance of death in a two-car collision.
Center of gravity
Some crash types tend to have more serious consequences. Rollovers have become more common in recent years, perhaps due to the increased popularity of taller SUVs, people carriers, and minivans, which have a higher center of gravity than standard passenger cars. Rollovers can be fatal, especially if the occupants are ejected because they were not wearing seat belts (83% of ejections during rollovers were fatal when the driver did not wear a seat belt, compared to 25% when they did).
After a first-generation Mercedes-Benz A-Class notoriously failed a 'moose test' (sudden swerving to avoid an obstacle) in 1997, some manufacturers enhanced suspension using stability control linked to an anti-lock braking system to reduce the likelihood of rollover. After retrofitting these systems to its models in 1999–2000, Mercedes saw its models involved in fewer crashes.
Now, about 40% of new US vehicles, mainly the SUVs, vans and pickup trucks that are more susceptible to rollover, are being produced with a lower center of gravity and enhanced suspension with stability control linked to its anti-lock braking system to reduce the risk of rollover and meet US federal requirements that mandate anti-rollover technology by September 2011.
Motorcycles
Motorcyclists and pillion-riders have little protection other than their clothing and helmets. This difference is reflected in the casualty statistics, where they are more than twice as likely to suffer severely after a collision. In 2005, there were 198,735 road crashes with 271,017 reported casualties on roads in Great Britain. This included 3,201 deaths (1.1%) and 28,954 serious injuries (10.7%) overall.
Of these casualties 178,302 (66%) were car users and 24,824 (9%) were motorcyclists, of whom 569 were killed (2.3%) and 5,939 seriously injured (24%).
Sociological factors
Studies in United States have shown that poor people have a greater risk of dying in a car crash than people who are well-off. Car deaths are also higher in poorer states.
Similar studies in France or Israel have shown the same results. This may be due to working-class people having less access to secure equipment in cars, having older cars which are less protected against crash, and needing to cover more distance to go to work each day.
COVID-19 lockdown impact
While the advent of the COVID lockdown meant a decrease in road traffic in the United States, the rates of incidents, speeding, and traffic fatalities rose in 2020 and 2021 (rate as measured against vehicle miles traveled). The traffic fatality rate jumped to 1.25 per 100 million vehicle miles traveled, up from 1.06 during the same period in 2019. Reasons cited for the increases are greater speeds, not wearing seatbelts, and driving while impaired.
In their preliminary report covering the first six months of 2021, the US nonprofit public safety advocacy group, the National Safety Council (NSC) estimated of total motor-vehicle deaths for the first six months of 2021 were 21,450, up 16% from 2020 and up 17% from 18,384 in 2019. The estimated mileage death rate in 2021 was 1.43 deaths per 100 million vehicle miles traveled, up 3% from 1.39 in 2020 and up 24% from 1.15 in 2019.
Preliminary data also show that even as traffic levels returned to normal after the onset of COVID in March–April 2020, drivers continued to drive at excessive speeds. A 2020 study conducted by INRIX, private company that analyzes traffic patters, behaviors and congestion, showed that as traffic levels returned to normal during the three-month period August to October 2020, growth in collisions (57%), outpaced the growth in miles traveled (22%) resulting in a higher than normal collision rate during this period.
In France, the Ministry of Interior reported that traffic incidents, crash-related injuries, and fatalities dropped in 2020 compared with 2019. Fatalities dropped 21.4%, injuries dropped 20.9%, and incidents overall dropped 20%. It also reported that the number of vehicles on the road dropped by 75%, which suggests the rate (incidents per vehicle-mile) increased.
Other
Other possibly hazardous factors that may alter a driver's soundness on the road include:
Irritability
Following specifically distinct rules too bureaucratically, inflexibly or rigidly when unique circumstances might suggest otherwise
Sudden swerving into somebody's blind spot without first clearly making oneself visible through the wing mirror
Unfamiliarity with one's dashboard features, center console or other interior handling devices after a recent car purchase
Lack of visibility due to windshield design, dense fog or sun glare
People-watching.
Traffic safety culture, a variety of aspects of safety culture could impact on the number of crashes.
Prevention
A large body of knowledge has been amassed on how to prevent car crashes, and reduce the severity of those that do occur.
United Nations
Owing to the global and massive scale of the issue, with predictions that by 2020 road traffic deaths and injuries will exceed HIV/AIDS as a cause of death and disability, the United Nations and its subsidiary bodies have passed resolutions and held conferences on the issue. The first United Nations General Assembly resolution and debate was in 2003 The World Day of Remembrance for Road Traffic Victims was declared in 2005. In 2009 the first high level ministerial conference on road safety was held in Moscow.
The World Health Organization, a specialized agency of the United Nations Organization, in its Global Status Report on Road Safety 2009, estimates that over 90% of the world's fatalities on the roads occur in low-income and middle-income countries, which have only 48% of the world's registered vehicles, and predicts road traffic injuries will rise to become the fifth leading cause of death by 2030.
The United Nations' Sustainable Development Goal 3, target 3.6 is directed at reducing road injuries and deaths. February 2020 saw a global ministerial conference which brought the Stockholm Declaration, setting a target to reduce global traffic deaths and injuries by 50% within ten years. The decade of 2021–2030 was declared the second decade of road safety.
Collision migration
Collisions migration refers to a situation where action to reduce road traffic collisions in one place may result in those collisions resurfacing elsewhere. For example, an accident blackspot may occur at a dangerous bend. The treatment for this may be to increase signage, post an advisory speed limit, apply a high-friction road surface, add crash barriers or any one of a number of other visible interventions. The immediate result may be to reduce collisions at the bend, but the subconscious relaxation on leaving the "dangerous" bend may cause drivers to act with fractionally less care on the rest of the road, resulting in an increase in collisions elsewhere on the road, and no overall improvement over the area. In the same way, increasing familiarity with the treated area will often result in a reduction over time to the previous level of care and may result in faster speeds around the bend due to perceived increased safety (risk compensation).
Epidemiology
In 2004 50 million more were injured in motor vehicle collisions. In 2013, between 1.25 million and 1.4 million people were killed in traffic collisions, up from 1.1 million deaths in 1990. That number represents about 2.5% of all deaths. Approximately 50 million additional people were injured in traffic collisions, a number unchanged from 2004.
India recorded 105,000 traffic deaths in a year, followed by China with over 96,000 deaths. This makes motor vehicle collisions the leading cause of injury and death among children worldwide 10–19 years old (260,000 children die a year, 10 million are injured) and the sixth leading preventable cause of death in the United States. In 2019, there were 36,096 people killed and 2.74 million people injured in motor vehicle traffic crashes on roadways in the United States. In the state of Texas alone, there were a total of 415,892 traffic collisions, including 3,005 fatal crashes in 2012. In Canada, they are the cause of 48% of severe injuries.
Crash rates
The safety performance of roadways is almost always reported as a rate. That is, some measure of harm (deaths, injuries, or number of crashes) divided by some measure of exposure to the risk of this harm. Rates are used so the safety performance of different locations can be compared, and to prioritize safety improvements.
Common rates related to road traffic fatalities include the number of deaths per capita, per registered vehicle, per licensed driver, or per vehicle mile or kilometer traveled. Simple counts are almost never used. The annual count of fatalities is a rate, namely, the number of fatalities per year.
There is no one rate that is superior to others in any general sense, it depends on the question asked and the available data. Some agencies concentrate on crashes per total vehicle distance traveled and others combine rates. Iowa, for example, selects high collision locations based on a combination of crashes per million miles traveled, crashes per mile per year, and value loss (crash severity).
Fatality
The definition of a road-traffic fatality varies from country to country. In the United States, the definition used in the Fatality Analysis Reporting System (FARS) run by the National Highway Traffic Safety Administration (NHTSA) is a person who dies within 30 days of a crash on a US public road involving a vehicle with an engine, the death being the result of the crash.
In the U.S., therefore, if a driver has a non-fatal heart attack that leads to a road-traffic crash that causes death, that is a road-traffic fatality. If the heart attack causes death prior to the crash, it is not a road-traffic fatality.
The definition of a road-traffic fatality can change with time in the same country. For example, fatality was defined in France as a person who dies in the six days (pre 2005) after the collision and was subsequently changed to the 30 days (post 2005) after the collision.
History
The world's first recorded road traffic death involving a motor vehicle occurred on 31 August 1869. Irish scientist Mary Ward died when she fell out of her cousins' steam car and was run over by it.
The British road engineer J. J. Leeming, compared the statistics for fatality rates in Great Britain, for transport-related incidents both before and after the introduction of the motor vehicle, for journeys, including those once by water that now are undertaken by motor vehicle: For the period 1863–1870 there were: 470 fatalities per million of population (76 on railways, 143 on roads, 251 on water); for the period 1891–1900 the corresponding figures were: 348 (63, 107, 178); for the period 1931–1938: 403 (22, 311, 70) and for the year 1963: 325 (10, 278, 37). Leeming concluded that the data showed that "travel accidents may even have been more frequent a century ago than they are now, at least for men".
He also compared the circumstances around road deaths as reported in various American states before the widespread introduction of speed limits and drunk-driving laws.
United States judges prioritized pedestrians' rights in city streets when early 20th century automobiles appeared. Pedestrian injuries were regarded as the fault of a motorist driving too fast. As automobile ownership increased, the rate of traffic deaths in the United States doubled from 1915 to 1921 when it reached 12 deaths per 100,000 Americans. The right to walk was considered dispensable a century later in 2021, when the annual death rate was 12.9 per 100,000. Safety focus on protecting the occupants of automobiles has victimized bicyclists and pedestrians whose injuries are attributed to individual carelessness. From 2010 to 2019, fatalities rose 36% for bicyclists and nearly doubled for those on foot. Reasons include larger vehicles, faster driving, and digital distractions making walking and biking in the United States far more dangerous than in other comparable nations.
The world's first autonomous car incident resulting in the death of a pedestrian occurred on 18 March 2018 in Arizona. The pedestrian was walking her bicycle outside of the crosswalk, and died in the hospital after she was struck by a self-driving car being tested by Uber.
Society and culture
Economic costs
The global economic cost of MVCs was estimated at $518 billion per year in 2003, and $100 billion in developing countries. The Centers for Disease Control and Prevention estimated the U.S. cost in 2000 at $230 billion. A 2010 US report estimated costs of $277 billion, which included lost productivity, medical costs, legal and court costs, emergency service costs (EMS), insurance administration costs, congestion costs, property damage, and workplace losses. "The value of societal harm from motor vehicle crashes, which includes both economic impacts and valuation for lost quality-of-life, was $870.8 billion in 2010. Sixty-eight percent of this value represents lost quality-of-life, while 32 percent are economic impacts."
Traffic collision affect the national economy as the cost of road injuries are estimated to account for 1.0% to 2.0% of the gross national product (GNP) of every country each year. A recent study from Nepal showed that the total economic costs of road injuries were approximately $122.88 million, equivalent to 1.52% of the total Nepal GNP for 2017, indicating the growing national financial burden associate with preventable road injuries and deaths.
The economic cost to the individuals involved in an MVC varies widely depending on geographic distribution, and varies largely on depth of accident insurance cover, and legislative policy. In the UK for example, a survey conducted using 500 post-accident insurance policy customers, showed an average individual financial loss of £1300.00. This is due in part to voluntary excesses that are common tactics used to reduce overall premium, and in part due to under valuation of vehicles. By contrast, Australian insurance policy holders are subject to an average financial loss of $950.00 AUD.
Legal consequences
There are a number of possible legal consequences for causing a traffic collision, including:
Traffic citations: drivers who are involved in a collision may receive one or more traffic citations for improper driving conduct such as speeding, failure to obey a traffic control device, or driving under the influence of drugs or alcohol. Convictions for traffic violations are usually penalized with fines, and for more severe offenses, the suspension or revocation of driving privileges.
Civil lawsuits: a driver who causes a traffic collision may be sued for damages resulting from the collision, including damages to property and injuries to other persons. Companies can be held liable if their employees cause motor vehicle crashes under a theory of vicarious liability. Other times, injured people can file a product liability lawsuit against a company that designed or distributed a dangerous vehicle or car part.
Criminal prosecution: More severe driving misconduct, including impaired driving, may result in criminal charges against the driver. In the event of a fatality, a charge of vehicular homicide is occasionally prosecuted, especially in cases involving alcohol. Convictions for alcohol offenses may result in the revocation or long term suspension of the driver's license, and sometimes jail time, mandatory drug or alcohol rehabilitation, or both.
Fraud
Sometimes, people may make false insurance claims or commit insurance fraud by staging collisions or jumping in front of moving cars.
United Kingdom
In the United Kingdom, the Pre-Action Protocol for Low Value Personal Injury Claims in Road Traffic Accidents from 31 July 2013, otherwise known as the RTA Protocol, the "upper limit" is £25,000 for an accident which occurred on or after 31 July 2013; the limit under a previous version of the protocol was £10,000 for an accident which had occurred on or after 30 April 2010 but before 31 July 2013.
United States
Motor vehicle crashes are the leading cause of death in the workplace in the United States accounting for 35 percent of all workplace fatalities. In the United States, individuals involved in motor vehicle collisions may be held financially liable for the consequences of a collision, including property damage, and injuries to passengers and drivers. Where another driver's vehicle is damaged as the result of a crash, some states allow the owner of the vehicle to recover both the cost of repair for the diminished value of the vehicle from the at-fault driver. Because the financial liability that results from causing a crash is so high, most U.S. states require drivers to carry liability insurance to cover these potential costs. In the event of serious injuries or fatalities, it is possible for injured persons to seek compensation in excess of the at-fault driver's insurance coverage.
Liability rules vary from state to state, with some laws adopting a tort system and others a no-fault insurance system. Most use a tort-based system, wherein injured people seek financial compensation from at-fault parties' insurance carriers. Twelve states take the no-fault approach, where injured parties file their primary claims with their own insurer.
Tort reform has changed the legal landscape. For example, Michigan had a unique no-fault system that guaranteed lifetime benefits for people injured in motor vehicle collisions. This changed in 2020, when the state legislature amended the laws, allowing people to opt for less coverage. While the aim of these laws is to reduce the cost of insurance premiums, catastrophically injured claimants might find themselves underinsured.
In some cases, involving a defect in the design or manufacture of motor vehicles, such as where defective design results in SUV rollovers or sudden unintended acceleration, crashes caused by defective tires, or where injuries are caused or worsened as a result of defective airbags, it is possible that the manufacturer will face a class action lawsuit.
Art
Cars have come to represent a part of the American Dream of ownership coupled with the freedom of the road. The violence of a car wreck provides a counterpoint to that promise and is the subject of artwork by a number of artists, such as John Salt and Li Yan. Though English, John Salt was drawn to American landscapes of wrecked vehicles like Desert Wreck (airbrushed oil on linen, 1972). Similarly, Jan Anders Nelson works with the wreck in its resting state in junkyards or forests, or as elements in his paintings and drawings. American Landscape is one example of Nelson's focus on the violence of the wreck with cars and trucks piled into a heap, left to the forces of nature and time. This recurring theme of violence is echoed in the work of Li Yan. His painting Accident Nº 6 looks at the energy released during a crash.
Andy Warhol used newspaper pictures of car wrecks with dead occupants in a number of his Disaster series of silkscreened canvases. John Chamberlain used components of wrecked cars (such as bumpers and crumpled sheet metal fenders) in his welded sculptures.
Crash is a 1973 novel by English author J. G. Ballard about car-crash sexual fetishism that was made into a film by David Cronenberg in 1996.
| Technology | Road transport | null |
23713759 | https://en.wikipedia.org/wiki/Hodgkin%20lymphoma | Hodgkin lymphoma | Hodgkin lymphoma (HL) is a type of lymphoma in which cancer originates from a specific type of white blood cell called lymphocytes, where multinucleated Reed–Sternberg cells (RS cells) are present in the patient's lymph nodes. The condition was named after the English physician Thomas Hodgkin, who first described it in 1832. Symptoms may include fever, night sweats, and weight loss. Often, nonpainful enlarged lymph nodes occur in the neck, under the arm, or in the groin. Persons affected may feel tired or be itchy.
The two major types of Hodgkin lymphoma are classic Hodgkin lymphoma and nodular lymphocyte-predominant Hodgkin lymphoma. About half of cases of Hodgkin lymphoma are due to Epstein–Barr virus (EBV) and these are generally the classic form. Other risk factors include a family history of the condition and having HIV/AIDS. Diagnosis is conducted by confirming the presence of cancer and identifying RS cells in lymph node biopsies. The virus-positive cases are classified as a form of the Epstein–Barr virus-associated lymphoproliferative diseases.
Hodgkin lymphoma may be treated with chemotherapy, radiation therapy, and stem-cell transplantation. The choice of treatment often depends on how advanced the cancer has become and whether or not it has favorable features. If the disease is detected early, a cure is often possible. In the United States, 88% of people diagnosed with Hodgkin lymphoma survive for five years or longer. For those under the age of 20, rates of survival are 97%. Radiation and some chemotherapy drugs, however, increase the risk of other cancers, heart disease, or lung disease over the subsequent decades.
In 2015, about 574,000 people globally had Hodgkin lymphoma, and 23,900 (4.2%) died. In the United States, 0.2% of people are affected at some point in their life. Most people are diagnosed with the disease between the ages of 20 and 40.
Signs and symptoms
People with Hodgkin lymphoma may present with these symptoms:
Lymphadenopathy: The most common symptom of Hodgkin is the painless enlargement of one or more lymph nodes. The nodes may also feel rubbery and swollen when examined. The nodes of the neck, armpits and groin (cervical and supraclavicular) are most frequently involved (80–90% of the time, on average). The lymph nodes of the chest are often affected, and these may be noticed on a chest radiograph.
Systemic symptoms: About one-third of people with Hodgkin disease may also present with systemic symptoms, including:
Itchy skin
Night sweats
Unexplained weight loss of at least 10% of the person's total body mass in six months or less
Low-grade fever.
Fatigue (lassitude)
Systemic symptoms such as fever, night sweats, and weight loss are known as B symptoms; thus, presence of these indicate that the person's stage is, for example, 2B instead of 2A.
Splenomegaly: Enlargement of the spleen is often present in people with Hodgkin lymphoma. The enlargement is seldom massive, and the size of the spleen may fluctuate during the course of treatment.
Hepatomegaly: Enlargement of the liver, due to liver involvement, is infrequent in people with Hodgkin lymphoma.
Hepatosplenomegaly: The enlargement of both the liver and spleen can be caused by the same disease.
Pain following alcohol consumption: Classically, involved nodes are painful after alcohol consumption, though this phenomenon is very uncommon, occurring in only two to three percent of people with Hodgkin lymphoma, thus having a low sensitivity. On the other hand, its positive predictive value is high enough for it to be regarded as a pathognomonic sign of Hodgkin lymphoma. The pain typically has an onset within minutes after ingesting alcohol, and is usually felt as coming from the vicinity where there is an involved lymph node. The pain has been described as either sharp and stabbing or dull and aching.
Back pain: Nonspecific back pain (pain that cannot be localised or its cause determined by examination or scanning techniques) has been reported in some cases of Hodgkin lymphoma. The lower back is most often affected.
Cyclical fever: People may also present with a cyclical high-grade fever known as the Pel–Ebstein fever, or more simply "P-E fever". However, there is debate as to whether the P-E fever truly exists.
Nephrotic syndrome can occur in individuals with Hodgkin lymphoma and is most commonly caused by minimal change disease.
May present with airway obstruction, pleural/pericardial effusion, hepatocellular dysfunction, or bone-marrow infiltration.
Diagnosis
Hodgkin lymphoma must be distinguished from noncancerous causes of lymph node swelling (such as various infections) and from other types of cancer. Definitive diagnosis is by lymph node biopsy (usually excisional biopsy with microscopic examination). Blood tests are also performed to assess function of major organs and safety for chemotherapy. Positron emission tomography (PET) is used to detect small deposits that do not show on CT scanning. PET scans are also useful in functional imaging (by using a radiolabeled glucose to image tissues of high metabolism). In some cases, a gallium scan may be used instead of a PET scan.
Types
The two main types of Hodgkin lymphoma are classic Hodgkin lymphoma and nodular lymphocyte-predominant Hodgkin lymphoma. The prevalence of classic Hodgkin lymphoma and nodular-lymphocyte Hodgkin lymphoma are about 90% and 10%, respectively. The morphology, phenotype, molecular features, and, therefore, the clinical behaviour and presentation of the two types differ.
Classic
Classic Hodgkin lymphoma (excluding nodular lymphocyte predominant Hodgkin lymphoma) can be subclassified into four pathologic subtypes based upon Reed–Sternberg cell morphology and the composition of the reactive cell infiltrate seen in the lymph node biopsy specimen (the cell composition around the Reed–Sternberg cell(s)). Presence of EBV in Reed-Sternberg cells is most commonly found in the subtypes lymphocyte depleted HL (>70%) and mixed cellularity HL (70%), whilst being less prevalent in lymphocyte-rich HL (40%) and relatively uncommon by comparison in nodular sclerosing HL.
For the other forms, although the traditional B-cell markers (such as CD20) are not expressed on all cells, Reed–Sternberg cells are usually of B cell origin. Although Hodgkin's is now frequently grouped with other B-cell malignancies, some T-cell markers (such as CD2 and CD4) are occasionally expressed. However, this may be an artifact of the ambiguity inherent in the diagnosis.
Hodgkin cells produce interleukin-21 (IL-21), which was once thought to be exclusive to T-cells. This feature may explain the behavior of classic Hodgkin lymphoma, including clusters of other immune cells gathered around HL cells (infiltrate) in cultures.
Nodular lymphocyte predominant
Nodular lymphocyte predominant Hodgkin lymphoma (NLPHL) is another subtype of Hodgkin lymphoma distinct from classic Hodgkin lymphoma and is characterized by the presence of popcorn cells which express CD20. Due to these differences, among others, NLPHL is often treated differently from classic Hodgkin lymphoma, including using rituximab in combination with AVBD chemotherapy, though individual cases vary and clinical trials are ongoing.
Staging
The staging is the same for both Hodgkin and non-Hodgkin lymphomas.
After Hodgkin lymphoma is diagnosed, a person will be staged: that is, they will undergo a series of tests and procedures that will determine what areas of the body are affected. These procedures may include documentation of their histology, a physical examination, blood tests, chest X-ray radiographs, computed tomography (CT)/positron emission tomography (PET)/magnetic resonance imaging (MRI) scans of the chest, abdomen and pelvis, and usually a bone marrow biopsy. PET scan is now used instead of the gallium scan for staging. On the PET scan, sites involved with lymphoma light up very brightly enabling accurate and reproducible imaging. In the past, a lymphangiogram or surgical laparotomy (which involves opening the abdominal cavity and visually inspecting for tumors) were performed. Lymphangiograms or laparotomies are very rarely performed, having been supplanted by improvements in imaging with the CT scan and PET scan.
On the basis of this staging, the person will be classified according to a staging classification (the Ann Arbor staging classification scheme is a common one):
Stage I is involvement of a single lymph node region (I) (mostly the cervical region) or single extralymphatic site (Ie);
Stage II is involvement of two or more lymph node regions on the same side of the diaphragm (II) or of one lymph node region and a contiguous extralymphatic site (IIe);
Stage III is involvement of lymph node regions on both sides of the diaphragm, which may include the spleen (IIIs) or limited contiguous extralymphatic organ or site (IIIe, IIIes);
Stage IV is disseminated involvement of one or more extralymphatic organs.
The absence of systemic symptoms is signified by adding "A" to the stage; the presence of systemic symptoms is signified by adding "B" to the stage. For localised extranodal extension from mass of nodes that does not advance the stage, subscript "E" is added. Splenic involvement is signified by adding "S" to the stage. The inclusion of "bulky disease" is signified by "X".
Pathology
Macroscopy
Affected lymph nodes (most often, laterocervical lymph nodes) are enlarged, but their shape is preserved because the capsule is not invaded. Usually, the cut surface is white-grey and uniform; in some histological subtypes (e.g. nodular sclerosis) a nodular aspect may appear.
A fibrin ring granuloma may be seen.
Microscopy
Microscopic examination of the lymph node biopsy reveals complete or partial effacement of the lymph node architecture by scattered large malignant cells known as Reed-Sternberg cells (RSC) (typical and variants) admixed within a reactive cell infiltrate composed of variable proportions of lymphocytes, histiocytes, eosinophils, and plasma cells. The Reed–Sternberg cells are identified as large often bi-nucleated cells with prominent nucleoli and an unusual CD45−, CD30+, CD15+/− immunophenotype. In approximately 50% of cases, the Reed–Sternberg cells are infected by the Epstein–Barr virus.
Characteristics of classic Reed–Sternberg cells include large size (20–50 micrometres), abundant, amphophilic, finely granular/homogeneous cytoplasm; two mirror-image nuclei (owl eyes) each with an eosinophilic nucleolus and a thick nuclear membrane (chromatin is distributed close to the nuclear membrane). Almost all of these cells have an increased copy number of chromosome 9p/9p24.1.
Variants:
Hodgkin cell (atypical mononuclear RSC) is a variant of RS cell, which has the same characteristics but is mononucleated.
Lacunar RSC is large, with a single hyperlobulated nucleus, multiple, small nucleoli and eosinophilic cytoplasm which is retracted around the nucleus, creating an empty space ("lacunae").
Pleomorphic RSC has multiple irregular nuclei.
"Popcorn" RSC (lympho-histiocytic variant) is a small cell, with a very lobulated nucleus, small nucleoli.
"Mummy" RSC has a compact nucleus with no nucleolus and basophilic cytoplasm.
Hodgkin lymphoma can be sub-classified by histological type. The cell histology in Hodgkin lymphoma is not as important as it is in non-Hodgkin lymphoma: the treatment and prognosis in classic Hodgkin lymphoma usually depends on the stage of disease rather than the histotype.
Management
The current approach for treatment aims to reduce the acute and long-term toxicities associated with Hodgkin lymphoma (e.g. cardiac damage and secondary cancers) and increase overall survival.
People with early stage disease (IA or IIA) can be treated effectively with radiation therapy or chemotherapy. The choice of treatment depends on the age, sex, bulk and the histological subtype of the disease.
Adding localised radiation therapy after the chemotherapy regimen may provide a longer progression-free survival compared with chemotherapy treatment alone. People with later disease (III, IVA, or IVB) are treated with combination chemotherapy alone. People of any stage with a large mass in the chest are usually treated with combined chemotherapy and radiation therapy.
The common non-Hodgkin treatment, rituximab (which is a monoclonal antibody against CD20) is not routinely used to treat Hodgkin lymphoma due to the lack of CD20 surface antigens in most cases. The use of rituximab in Hodgkin lymphoma, including the lymphocyte predominant subtype has been recently reviewed. The evidence is very uncertain about the effect of Nivolumab for patients with a Hodgkin's lymphoma e.g. on the overall survival.
Increased age is an adverse risk factor for Hodgkin lymphoma, but in general elderly people (≥ 60 years of age) without major comorbidities are sufficiently fit to tolerate therapy with curative intent. Despite this, treatment outcome in the elderly patient is not comparable to that of younger people and the disease is a different entity in older people where different considerations enter into treatment decisions.
Recently, two novel targeted drugs have been developed for relapsing and refractory HL patients; Brentuximab vedotin, a CD30 antibody conjugated with a cytotoxic component MMAE, and the checkpoint inhibitors, Nivolumab and Pembrolizumab. This has been an important step in the treatment for the few, but still existing refractory patients.
For Hodgkin lymphomas, radiation oncologists typically use external beam radiation therapy (sometimes shortened to EBRT or XRT). Radiation oncologists deliver external beam radiation therapy to the lymphoma from a machine called a linear accelerator which produces high energy X-rays and electrons. People usually describe treatments as painless and similar to getting an X-ray. Treatments last less than 30 minutes each.
For lymphomas, there are a few different ways radiation oncologists target the cancer cells. Involved site radiation is when the radiation oncologists give radiation only to those parts of the person's body known to have the cancer. Very often, this is combined with chemotherapy. Radiation therapy directed above the diaphragm to the neck, chest or underarms is called mantle field radiation. Radiation to below the diaphragm to the abdomen, spleen or pelvis is called inverted-Y field radiation. Total nodal irradiation is when the therapist gives radiation to all the lymph nodes in the body to destroy cells that may have spread.
Adverse effects
The high cure rates and long survival of many people with Hodgkin lymphoma has led to a high concern with late adverse effects of treatment, including cardiovascular disease and second malignancies such as acute leukemias, lymphomas, and solid tumors within the radiation therapy field. Most people with early-stage disease are now treated with abbreviated chemotherapy and involved site radiation therapy rather than with radiation therapy alone. Clinical research strategies are exploring reduction of the duration of chemotherapy and dose and volume of radiation therapy in an attempt to reduce late morbidity and mortality of treatment while maintaining high cure rates. Hospitals are also treating those who respond quickly to chemotherapy with no radiation.
In childhood cases of Hodgkin lymphoma, long-term endocrine adverse effects are a major concern, mainly gonadal dysfunction and growth retardation. Gonadal dysfunction seems to be the most severe endocrine long-term effect, especially after treatment with alkylating agents or pelvic radiotherapy.
It is possible that patients undergoing a chemotherapy need a platelet transfusion. If a stem cell transplantation is necessary for the treatment of a relapse, graft-versus-host diseases might occur.
Supportive treatment
Adding physical exercises to the standard treatment for adult patients with haematological malignancies like Hodgkin lymphoma may result in little to no difference in the mortality, the quality of life and the physical functioning. These exercises may result in a slight reduction in depression. Furthermore, aerobic physical exercises probably reduce fatigue. The evidence is very uncertain about the effect on anxiety and serious adverse events.
Prognosis
Treatment of Hodgkin's disease has been improving over the past few decades. Recent trials that have made use of new types of chemotherapy have indicated higher survival rates than have previously been seen. In a 2007 European trial, the five-year survival rate for those people with a favorable prognosis (FFP) was 98%, while that for people with worse outlooks was at least 85%.
In 1998, an international effort identified seven prognostic factors that accurately predict the success rate of conventional treatment in people with locally extensive or advanced-stage Hodgkin lymphoma. Freedom from progression (FFP) at five years was directly related to the number of factors present in a person. The five-year FFP for people with zero factors is 84%. Each additional factor lowers the five-year FFP rate by 7%, such that the five-year FFP for a person with five or more factors is 42%.
The adverse prognostic factors identified in the international study are:
Age ≥ 45 years
Stage IV disease
Hemoglobin < 10.5 g/dl
Lymphocyte count < 600/μL or < 8%
Male
Albumin < 4.0 g/dl
White blood count ≥ 15,000/μL
Other studies have reported the following to be the most important adverse prognostic factors: mixed-cellularity or lymphocyte-depleted histologies, male sex, large number of involved nodal sites, advanced stage, age of 40 years or more, the presence of B symptoms, high erythrocyte sedimentation rate, and bulky disease (widening of the mediastinum by more than one third, or the presence of a nodal mass measuring more than 10 cm in any dimension.)
More recently, the use of positron emission tomography (PET) early after commencing chemotherapy has demonstrated to have powerful prognostic ability. This enables assessment of an individual's response to chemotherapy as the PET activity switches off rapidly in people who are responding. In this study, after two cycles of ABVD chemotherapy, 83% of people were free of disease at 3 years if they had a negative PET versus only 28% in those with positive PET scans. This prognostic method improves on FFP estimates based on the seven conventional factors. Several trials are underway to see if PET-based risk adapted response can be used to improve a person's outcomes by changing chemotherapy early in people who are not responding.
The evidence is very uncertain about the effect of negative (= good prognosis) or positive (= bad prognosis) interim PET scan results for patients with a Hodgkin's lymphoma on the progression-free survival. Negative interim PET scan results may result in an increase in progression-free survival compared if the adjusted result was measured. Negative interim PET scan results probably result in a large increase in the overall survival compared to those with a positive interim PET scan result,
Epidemiology
Unlike some other lymphomas, whose number of new cases per year increases with age, Hodgkin lymphoma has a bimodal curve for the number of cases; that is, it occurs most frequently in two separate age groups, the first being young adulthood (age 15–35) and the second being in those over 55 years old although these peaks may vary slightly with nationality. Overall, it is more common in males, except for the nodular sclerosis variant, which is slightly more common in females. The annual number of cases of Hodgkin lymphoma is 2.7 per 100,000 per persons per year, and the disease accounts for slightly less than 1% of all cancers worldwide.
In 2010, globally it resulted in about 18,000 deaths down from 19,000 in 1990. In 2012, there were an estimated 65,950 cases and 25,469 deaths from Hodgkin lymphoma worldwide, with 28,852 and 37,098 cases occurring in developed and developing countries, respectively. However, the age-standardized rates were higher in developed regions, with the greatest rates in the Americas (1.5 per 100,000), East Mediterranean Region (1.5 per 100,000), and Europe (2.0 per 100,000). The East Mediterranean Region also has the highest age-standardized mortality rate of 1.0 per 100,000, which is mainly attributed to lifestyle and environmental risk factors associated with transitional economies such as smoking, obesity, physical inactivity, and reproductive behaviors, as well as availability of diagnostic practices and awareness of the disease.
The number of cases of Hodgkin lymphoma is increased in people with HIV infection. In contrast to many other lymphomas associated with HIV infection it occurs most commonly in people with higher CD4 T cell counts.
Canada
Hodgkin lymphoma accounts for 0.6% of all male cancer cases, and 0.4% of all female cancer cases in Canada. In 2017, approximately 990 Canadians will be diagnosed with Hodgkin lymphoma, and 140 will die of the disease.
UK
Hodgkin lymphoma accounts for less than 1% of all cancer cases and deaths in the UK. Around 1,800 people were diagnosed with the disease in 2011, and around 330 people died in 2012.
United States
In 2016, there were 8,389 new cases and 1,000 mortalities attributed to Hodgkin lymphoma, a decrease from the 8,625 new cases and 1,120 mortalities in 2015. As of January 1, 2016, the 5-year limited duration prevalence of Hodgkin lymphoma was 37,513 representing 0.71% of all diagnosed cancers in the U.S.
History
Hodgkin lymphoma was first described in an 1832 report by Thomas Hodgkin, although Hodgkin noted that perhaps an earlier reference to the condition was provided by Marcello Malpighi in 1666. While occupied as museum curator at Guy's Hospital, London, Hodgkin studied seven people with painless lymph node enlargement. Of the seven cases, two were under the care of Richard Bright, one was of Thomas Addison, and one was of Robert Carswell. Carswell's report of the seventh case was accompanied by numerous illustrations that aided early descriptions of the disease.
Hodgkin's report on the seven cases, entitled "On some morbid appearances of the absorbent glands and spleen", was presented to the Medical and Chirurgical Society of London in January 1832 and was subsequently published in the society's journal, Medical-Chirurgical Society Transactions. Hodgkin's paper went largely unnoticed, however, even though Bright highlighted it in an 1838 publication. Indeed, Hodgkin himself did not view his contribution as particularly significant.
In 1856, Samuel Wilks independently reported on a series of patients with the same disease that Hodgkin had previously described. Wilks, a successor to Hodgkin at Guy's Hospital, was unaware of Hodgkin's prior work on the subject. Bright informed Wilks of Hodgkin's contribution and in 1865, Wilks published a second paper, entitled "Cases of enlargement of the lymphatic glands and spleen", in which he named the illness "Hodgkin's disease" in honor of his predecessor.
Theodor Langhans and WS Greenfield first described the microscopic characteristics of Hodgkin lymphoma in 1872 and 1878, respectively. In 1898 and 1902, respectively, Carl Sternberg and Dorothy Reed independently described the cytogenetic features of the malignant cells of Hodgkin lymphoma, now called Reed–Sternberg cells.
Tissue specimens from Hodgkin's seven cases were preserved at Guy's Hospital. Nearly 100 years after Hodgkin's initial publication, histopathologic reexamination confirmed Hodgkin lymphoma in only three of seven of these people. The remaining cases included non-Hodgkin lymphoma, tuberculosis, and syphilis.
Hodgkin lymphoma was one of the first cancers to be treated successfully with radiation therapy and, later, it was one of the first to be treated by combination chemotherapy.
Notable cases
Paul Allen, co-founder of Microsoft was diagnosed with Hodgkin lymphoma in 1982. He later died from non-Hodgkin lymphoma, on October 15, 2018.
Eric Berry, All-Pro strong safety for the Kansas City Chiefs of the National Football League, diagnosed in 2014.
David Brooks, Welsh professional footballer, diagnosed in 2021 while playing for AFC Bournemouth.
Dale Carnegie, Public Speaker and author of "How to Win Friends and Influence People" (c)1936, Born 1888, Maryville, MO, died 1955, Forest Hills,NY.
Howard Carter, Egyptologist and discoverer of the Tomb of Tutankhamun, died in 1939 from Hodgkin's disease.
Starchild Abraham Cherrix, a teenager whose refusal to undergo further conventional treatment after relapsing in 2006 resulted in a court battle and a change to Virginia laws about medical neglect.
James Conner, running back and 2014 ACC Player of the Year for the Arizona Cardinals.
Michael Cuccione, Canadian child actor, was diagnosed in 1994 at age 9. Treatments that rendered him cancer-free, including chemotherapy, a bone marrow transplant, and radiation, left him with permanent lung and respiratory problems and he died in 2001 just after turning 16.
Victoria Duval, American tennis player, was diagnosed in 2014.
Gerald Finzi, British composer, was diagnosed in 1951 and died in 1956.
Mist Edvardsdóttir, Icelandic football player and member of the Icelandic women's national team. Diagnosed in June 2014 at the age of 23. Continued to play until becoming too ill due to chemotherapy. Made recovery in early 2015.
Delta Goodrem, Australian singer, songwriter, and actress. She was diagnosed in July 2003 at the age of 18.
Hank Green, one of the cofounders of the Vlogbrothers, VidCon and production company Complexly, announced that he was diagnosed with Hodgkin lymphoma in a video he released on 19 May 2023. Green would announce his entrance into remission on 25 August 2023.
Jiří Grossmann, Czechoslovak theatre actor, poet, and composer
Michael C. Hall (born February 1, 1971), American actor, best known for his lead role as Dexter Morgan, in Showtime's crime series Dexter. In 2010, aged 38, Hall announced he was undergoing treatment for Hodgkin lymphoma; within two years, the disease was in full remission.
Richard Harris, Irish actor who portrayed Albus Dumbledore in the first two Harry Potter movies, died on October 25, 2002, after being diagnosed earlier that year.
Daniel Hauser, whose mother fled with him in 2009 in order to prevent him from undergoing chemotherapy.
Tessa James, Australian actress, was diagnosed in 2014.
Sean Kent, American stand up comedian and actor. Was diagnosed in 2002 while writing on The Best Damn Sports Show Period. After three months of chemotherapy and one month of radiation, the cancer went into remission.
Mario Lemieux, Hall of Fame NHL player, co-owner of the Pittsburgh Penguins and founder of the Mario Lemieux Foundation, diagnosed in 1993.
Dinu Lipatti (1917–1950), Romanian classical pianist and composer. Diagnosed in 1947, received cortisone treatment in 1949; died from a burst abscess on his one lung.
Jack Lisowski, English snooker player, diagnosed in 2008 at the age of 16.
Mamta Mohandas, Indian film actress and producer, diagnosed in 2010.
Nanni Moretti, Italian actor and director.
Laura Packard, health care activist diagnosed in 2017, spoke at the 2020 Democratic National Convention.
Nikola Pokrivač, Croatian soccer midfielder, diagnosed in 2015.
Anthony Rizzo, MLB All-Star first baseman for the New York Yankees, diagnosed in May 2008 while signed as a minor league player for the Boston Red Sox.
Dave Roberts, MLB outfielder and manager of the Los Angeles Dodgers. Diagnosed in March 2010 while he was a coach for the San Diego Padres.
Chip Roy, Texas congressman.
Flip Saunders, head coach of the NBA team Minnesota Timberwolves, announced in August 2015 that he was diagnosed with Hodgkin's disease. He died of the disease in October 2015.
Arlen Specter, United States Senator from Pennsylvania (1981–2011), diagnosed in 2005. He later died from non-Hodgkin lymphoma in 2012.
Brandon Tartikoff, American television executive, diagnosed around 1974, died in 1997.
Bernardo Tengarrinha, Portuguese professional footballer, diagnosed in 2017 Tengarrinha died on October 30, 2021, at the age of 32. Hours later, his former teams FC Porto and Boavista FC paid tribute to him before the local derby. playing for FC Porto.
Ethan Zohn, American professional soccer player and a winner of the Survivor reality television series. Zohn was diagnosed twice (in 2009 and 2011).
Richard Holliday, American professional wrestler, reported to be ill with diagnosis from June 2022, completed treatment February 17, 2023
Rick Czaplewski, Author Better Dirty Than Done, Motivational Speaker, 2009 Boston Marathon Finisher.
| Biology and health sciences | Cancer | Health |
6277878 | https://en.wikipedia.org/wiki/Open%20science | Open science | Open science is the movement to make scientific research (including publications, data, physical samples, and software) and its dissemination accessible to all levels of society, amateur or professional. Open science is transparent and accessible knowledge that is shared and developed through collaborative networks. It encompasses practices such as publishing open research, campaigning for open access, encouraging scientists to practice open-notebook science (such as openly sharing data and code), broader dissemination and engagement in science and generally making it easier to publish, access and communicate scientific knowledge.
Usage of the term varies substantially across disciplines, with a notable prevalence in the STEM disciplines. Open research is often used quasi-synonymously to address the gap that the denotion of "science" might have regarding an inclusion of the Arts, Humanities and Social Sciences. The primary focus connecting all disciplines is the widespread uptake of new technologies and tools, and the underlying ecology of the production, dissemination and reception of knowledge from a research-based point-of-view.
As Tennant et al. (2020) note, the term open science "implicitly seems only to regard ‘scientific’ disciplines, whereas open scholarship can be considered to include research from the Arts and Humanities, as well as the different roles and practices that researchers perform as educators and communicators, and an underlying open philosophy of sharing knowledge beyond research communities."
Open science can be seen as a continuation of, rather than a revolution in, practices begun in the 17th century with the advent of the academic journal, when the societal demand for access to scientific knowledge reached a point at which it became necessary for groups of scientists to share resources with each other. In modern times there is debate about the extent to which scientific information should be shared. The conflict that led to the Open Science movement is between the desire of scientists to have access to shared resources versus the desire of individual entities to profit when other entities partake of their resources. Additionally, the status of open access and resources that are available for its promotion are likely to differ from one field of academic inquiry to another.
Principles
The six principles of open science are:
Open methodology
Open source
Open data
Open access
Open peer review
Open educational resources
Background
Science is broadly understood as collecting, analyzing, publishing, reanalyzing, criticizing, and reusing data. Proponents of open science identify a number of barriers that impede or dissuade the broad dissemination of scientific data.
These include financial paywalls of for-profit research publishers, restrictions on usage applied by publishers of data, poor formatting of data or use of proprietary software that makes it difficult to re-purpose, and cultural reluctance to publish data for fears of losing control of how the information is used.
According to the FOSTER taxonomy Open science can often include aspects of Open access, Open data and the open source movement whereby modern science requires software to process data and information.
Open research computation also addresses the problem of reproducibility of scientific results.
Types
The term "open science" does not have any one fixed definition or operationalization. On the one hand, it has been referred to as a "puzzling phenomenon". On the other hand, the term has been used to encapsulate a series of principles that aim to foster scientific growth and its complementary access to the public. Two influential sociologists, Benedikt Fecher and Sascha Friesike, have created multiple "schools of thought" that describe the different interpretations of the term.
According to Fecher and Friesike ‘Open Science’ is an umbrella term for various assumptions about the development and dissemination of knowledge. To show the term's multitudinous perceptions, they differentiate between five Open Science schools of thought:
Infrastructure School
The infrastructure school is founded on the assumption that "efficient" research depends on the availability of tools and applications. Therefore, the "goal" of the school is to promote the creation of openly available platforms, tools, and services for scientists. Hence, the infrastructure school is concerned with the technical infrastructure that promotes the development of emerging and developing research practices through the use of the internet, including the use of software and applications, in addition to conventional computing networks. In that sense, the infrastructure school regards open science as a technological challenge. The infrastructure school is tied closely with the notion of "cyberscience", which describes the trend of applying information and communication technologies to scientific research, which has led to an amicable development of the infrastructure school. Specific elements of this prosperity include increasing collaboration and interaction between scientists, as well as the development of "open-source science" practices. The sociologists discuss two central trends in the infrastructure school:
1. Distributed computing: This trend encapsulates practices that outsource complex, process-heavy scientific computing to a network of volunteer computers around the world. The examples that the sociologists cite in their paper is that of the Open Science Grid, which enables the development of large-scale projects that require high-volume data management and processing, which is accomplished through a distributed computer network. Moreover, the grid provides the necessary tools that the scientists can use to facilitate this process.
2. Social and Collaboration Networks of Scientists: This trend encapsulates the development of software that makes interaction with other researchers and scientific collaborations much easier than traditional, non-digital practices. Specifically, the trend is focused on implementing newer Web 2.0 tools to facilitate research related activities on the internet. De Roure and colleagues (2008) list a series of four key capabilities which they believe define a Social Virtual Research Environment (SVRE):
The SVRE should primarily aid the management and sharing of research objects. The authors define these to be a variety of digital commodities that are used repeatedly by researchers.
Second, the SVRE should have inbuilt incentives for researchers to make their research objects available on the online platform.
Third, the SVRE should be "open" as well as "extensible", implying that different types of digital artifacts composing the SVRE can be easily integrated.
Fourth, the authors propose that the SVRE is more than a simple storage tool for research information. Instead, the researchers propose that the platform should be "actionable". That is, the platform should be built in such a way that research objects can be used in the conduct of research as opposed to simply being stored.
Measurement school
The measurement school, in the view of the authors, deals with developing alternative methods to determine scientific impact. This school acknowledges that measurements of scientific impact are crucial to a researcher's reputation, funding opportunities, and career development. Hence, the authors argue, that any discourse about Open Science is pivoted around developing a robust measure of scientific impact in the digital age. The authors then discuss other research indicating support for the measurement school. The three key currents of previous literature discussed by the authors are:
The peer-review is described as being time-consuming.
The impact of an article, tied to the name of the authors of the article, is related more to the circulation of the journal rather than the overall quality of the article itself.
New publishing formats that are closely aligned with the philosophy of Open Science are rarely found in the format of a journal that allows for the assignment of the impact factor.
Hence, this school argues that there are faster impact measurement technologies that can account for a range of publication types as well as social media web coverage of a scientific contribution to arrive at a complete evaluation of how impactful the science contribution was. The gist of the argument for this school is that hidden uses like reading, bookmarking, sharing, discussing and rating are traceable activities, and these traces can and should be used to develop a newer measure of scientific impact. The umbrella jargon for this new type of impact measurements is called altmetrics, coined in a 2011 article by Priem et al., (2011). Markedly, the authors discuss evidence that altmetrics differ from traditional webometrics which are slow and unstructured. Altmetrics are proposed to rely upon a greater set of measures that account for tweets, blogs, discussions, and bookmarks. The authors claim that the existing literature has often proposed that altmetrics should also encapsulate the scientific process, and measure the process of research and collaboration to create an overall metric. However, the authors are explicit in their assessment that few papers offer methodological details as to how to accomplish this. The authors use this and the general dearth of evidence to conclude that research in the area of altmetrics is still in its infancy.
Public School
According to the authors, the central concern of the school is to make science accessible to a wider audience. The inherent assumption of this school, as described by the authors, is that the newer communication technologies such as Web 2.0 allow scientists to open up the research process and also allow scientist to better prepare their "products of research" for interested non-experts. Hence, the school is characterized by two broad streams: one argues for the access of the research process to the masses, whereas the other argues for increased access to the scientific product to the public.
Accessibility to the Research Process: Communication technology allows not only for the constant documentation of research but also promotes the inclusion of many different external individuals in the process itself. The authors cite citizen science – the participation of non-scientists and amateurs in research. The authors discuss instances in which gaming tools allow scientists to harness the brain power of a volunteer workforce to run through several permutations of protein-folded structures. This allows for scientists to eliminate many more plausible protein structures while also "enriching" the citizens about science. The authors also discuss a common criticism of this approach: the amateur nature of the participants threatens to pervade the scientific rigor of experimentation.
Comprehensibility of the Research Result: This stream of research concerns itself with making research understandable for a wider audience. The authors describe a host of authors that promote the use of specific tools for scientific communication, such as microblogging services, to direct users to relevant literature. The authors claim that this school proposes that it is the obligation of every researcher to make their research accessible to the public. The authors then proceed to discuss if there is an emerging market for brokers and mediators of knowledge that is otherwise too complicated for the public to grasp.
Democratic school
The democratic school concerns itself with the concept of access to knowledge. As opposed to focusing on the accessibility of research and its understandability, advocates of this school focus on the access of products of research to the public. The central concern of the school is with the legal and other obstacles that hinder the access of research publications and scientific data to the public. Proponents assert that any research product should be freely available. and that everyone has the same, equal right of access to knowledge, especially in the instances of state-funded experiments and data. Two central currents characterize this school: Open Access and Open Data.
Open Data: Opposition to the notion that publishing journals should claim copyright over experimental data, which prevents the re-use of data and therefore lowers the overall efficiency of science in general. The claim is that journals have no use of the experimental data and that allowing other researchers to use this data will be fruitful. Only a quarter of researchers agree to share their data with other researchers because of the effort required for compliance.
Open Access to Research Publication: According to this school, there is a gap between the creation and sharing of knowledge. Proponents argue that even though scientific knowledge doubles every 5 years, access to this knowledge remains limited. These proponents consider access to knowledge as a necessity for human development, especially in the economic sense.
Pragmatic School
The pragmatic school considers Open Science as the possibility to make knowledge creation and dissemination more efficient by increasing the collaboration throughout the research process. Proponents argue that science could be optimized by modularizing the process and opening up the scientific value chain. 'Open' in this sense follows very much the concept of open innovation. Take for instance transfers the outside-in (including external knowledge in the production process) and inside-out (spillovers from the formerly closed production process) principles to science. Web 2.0 is considered a set of helpful tools that can foster collaboration (sometimes also referred to as Science 2.0). Further, citizen science is seen as a form of collaboration that includes knowledge and information from non-scientists. Fecher and Friesike describe data sharing as an example of the pragmatic school as it enables researchers to use other researchers' data to pursue new research questions or to conduct data-driven replications.
History
The widespread adoption of the institution of the scientific journal marks the beginning of the modern concept of open science. Before this time societies pressured scientists into secretive behaviors.
Before journals
Before the advent of scientific journals, scientists had little to gain and much to lose by publicizing scientific discoveries. Many scientists, including Galileo, Kepler, Isaac Newton, Christiaan Huygens, and Robert Hooke, made claim to their discoveries by describing them in papers coded in anagrams or cyphers and then distributing the coded text. Their intent was to develop their discovery into something off which they could profit, then reveal their discovery to prove ownership when they were prepared to make a claim on it.
The system of not publicizing discoveries caused problems because discoveries were not shared quickly and because it sometimes was difficult for the discoverer to prove priority. Newton and Gottfried Leibniz both claimed priority in discovering calculus. Newton said that he wrote about calculus in the 1660s and 1670s, but did not publish until 1693. Leibniz published "Nova Methodus pro Maximis et Minimis", a treatise on calculus, in 1684. Debates over priority are inherent in systems where science is not published openly, and this was problematic for scientists who wanted to benefit from priority.
These cases are representative of a system of aristocratic patronage in which scientists received funding to develop either immediately useful things or to entertain. In this sense, funding of science gave prestige to the patron in the same way that funding of artists, writers, architects, and philosophers did. Because of this, scientists were under pressure to satisfy the desires of their patrons, and discouraged from being open with research which would bring prestige to persons other than their patrons.
Emergence of academies and journals
Eventually the individual patronage system ceased to provide the scientific output which society began to demand. Single patrons could not sufficiently fund scientists, who had unstable careers and needed consistent funding. The development which changed this was a trend to pool research by multiple scientists into an academy funded by multiple patrons. In 1660 England established the Royal Society and in 1666 the French established the French Academy of Sciences. Between the 1660s and 1793, governments gave official recognition to 70 other scientific organizations modeled after those two academies. In 1665, Henry Oldenburg became the editor of Philosophical Transactions of the Royal Society, the first academic journal devoted to science, and the foundation for the growth of scientific publishing. By 1699 there were 30 scientific journals; by 1790 there were 1052. Since then publishing has expanded at even greater rates.
Popular Science Writing
The first popular science periodical of its kind was published in 1872, under a suggestive name that is still a modern portal for the offering science journalism: Popular Science. The magazine claims to have documented the invention of the telephone, the phonograph, the electric light and the onset of automobile technology. The magazine goes so far as to claim that the "history of Popular Science is a true reflection of humankind's progress over the past 129+ years". Discussions of popular science writing most often contend their arguments around some type of "Science Boom". A recent historiographic account of popular science traces mentions of the term "science boom" to Daniel Greenberg's Science and Government Reports in 1979 which posited that "Scientific magazines are bursting out all over. Similarly, this account discusses the publication Time, and its cover story of Carl Sagan in 1980 as propagating the claim that popular science has "turned into enthusiasm". Crucially, this secondary account asks the important question as to what was considered as popular "science" to begin with. The paper claims that any account of how popular science writing bridged the gap between the informed masses and the expert scientists must first consider who was considered a scientist to begin with.
Collaboration among academies
In modern times many academies have pressured researchers at publicly funded universities and research institutions to engage in a mix of sharing research and making some technological developments proprietary. Some research products have the potential to generate commercial revenue, and in hope of capitalizing on these products, many research institutions withhold information and technology which otherwise would lead to overall scientific advancement if other research institutions had access to these resources. It is difficult to predict the potential payouts of technology or to assess the costs of withholding it, but there is general agreement that the benefit to any single institution of holding technology is not as great as the cost of withholding it from all other research institutions.
Coining of term "Open Science"
Steve Mann claimed to have coined the term "Open Science" in 1998. He also registered the domain names openscience.com and openscience.org in 1998, which he sold to degruyter.com in 2011. The term was previously used in a manner that refers to today's 'open science' norms by Daryl E. Chubin in his 1985 essay "Open Science and Closed Science: Tradeoffs in a Democracy". Chubin's essay cited Robert K. Merton's 1942 proposal of what we now refer to as Mertonian Norms for ideal science practices and scientific modes of communication. The term was used sporadically in the 1970s and 1980s in various scholarship to refer to different things.
Internet and the free access to scientific documents
The open science movement, as presented in activist and institutional discourses at the beginning of the 21st century, refers to different ways of opening up science, especially in the Internet age. Its first pillar is free access to scientific publications. The Budapest conference organised by the Open Society Foundations in 2001 was decisive in imposing this issue on the political landscape. The resulting declaration calls for the use of digital tools such as open archives and open access journals, free of charge for the reader.
The idea of open access to scientific publications quickly became inseparable from the question of free licenses to guarantee the right to disseminate and possibly modify shared documents, such as the Creative Commons licenses, created in 2002. In 2011, a new text from the Budapest Open Initiative explicitly refers to the relevance of the CC-BY license to guarantee free dissemination and not only free access to a scientific document.
The openness promise by the Internet is then extended to research data, which underpins scientific studies in different disciplines, as mentioned already in the Berlin Declaration in 2003. In 2007, the Organisation for Economic Co-operation and Development (OECD) published a report on access to publicly funded research data, in which it defined it as the data that validates research results.
Beyond its democratic virtues, open science aims to respond to the replication crisis of research results, notably through the generalization of the opening of data or source code used to produce them or through the dissemination of methodological articles.
The open science movement inspired several regulatory and legislative measures. Thus, in 2007, the University of Liège made the deposit of its researchers’ publications in its institutional open repository (Orbi) compulsory. The next year, the NIH Public Access Policy adopted a similar mandate for every paper funded by the National Institutes of Health. In France, the law for a digital Republic enacted in 2016 creates the right to deposit the validated manuscript of a scientific article in an open archive, with an embargo period following the date of publication in the journal. The law also creates the principle of reuse of public data by default.
Politics
In many countries, governments fund some science research. Scientists often publish the results of their research by writing articles and donating them to be published in scholarly journals, which frequently are commercial. Public entities such as universities and libraries subscribe to these journals. Michael Eisen, a founder of the Public Library of Science, has described this system by saying that "taxpayers who already paid for the research would have to pay again to read the results."
In December 2011, some United States legislators introduced a bill called the Research Works Act, which would prohibit federal agencies from issuing grants with any provision requiring that articles reporting on taxpayer-funded research be published for free to the public online. Darrell Issa, a co-sponsor of the bill, explained the bill by saying that "Publicly funded research is and must continue to be absolutely available to the public. We must also protect the value added to publicly funded research by the private sector and ensure that there is still an active commercial and non-profit research community." One response to this bill was protests from various researchers; among them was a boycott of commercial publisher Elsevier called The Cost of Knowledge.
The Dutch Presidency of the Council of the European Union called out for action in April 2016 to migrate European Commission funded research to Open Science. European Commissioner Carlos Moedas introduced the Open Science Cloud at the Open Science Conference in Amsterdam on 4–5 April. During this meeting also The Amsterdam Call for Action on Open Science was presented, a living document outlining concrete actions for the European Community to move to Open Science. The European Commission continues to be committed to an Open Science policy including developing a repository for research digital objects, European Open Science Cloud (EOSC) and metrics for evaluating quality and impact.
In October2021, the French Ministry of Higher Education, Research and Innovation released an official translation of its second plan for open science spanning the years 2021–2024.
Standard setting instruments
There is currently no global normative framework covering all aspects of Open Science. In November 2019, UNESCO was tasked by its 193 Member States, during their 40th General Conference, with leading a global dialogue on Open Science to identify globally-agreed norms and to create a standard-setting instrument. The multistakeholder, consultative, inclusive and participatory process to define a new global normative instrument on Open Science is expected to take two years and to lead to the adoption of a UNESCO Recommendation on Open Science by Member States in 2021.
Two UN frameworks set out some common global standards for application of Open Science and closely related concepts: the UNESCO Recommendation on Science and Scientific Researchers, approved by the General Conference at its 39th session in 2017, and the UNESCO Strategy on Open Access to scientific information and research, approved by the General Conference at its 36th session in 2011.
Open Science and Research Assessment
A central aspect of the Open Science movement is the reform of research assessment. Initiatives such as the Coalition for Advancing Research Assessment (CoARA) and the San Francisco Declaration on Research Assessment (DORA) advocate moving away from traditional quantitative metrics like the Journal Impact Factor (JIF) and the h-Index, as these often exhibit biases and neglect qualitative aspects. Instead, alternative metrics and indicators, such as altmetrics and Open Science indicators, are to be given greater consideration. Open Science indicators include metrics such as the number of open access publications, data management plans, preprints, FAIR-licensed data, and open peer review reports. These approaches aim to promote the transparency and reusability of scientific outcomes, thereby enabling a fairer and more comprehensive evaluation of scientific achievements.While Open Science aims to enhance transparency, accessibility, and collaboration, the introduction of numerous new metrics to measure openness has led to unintended consequences. These metrics often rely on quantitative indicators, which conflict with the holistic and qualitative approaches advocated by initiatives such as CoARA and DORA. The core issue is that these metrics are designed not only to measure but also to influence researchers' behavior. This can result in "metric-driven" practices that undermine research quality. Additionally, Open Science metrics lack standardization and clarity regarding what they truly aim to measure. The risk is that while these metrics may incentivize openness, they could simultaneously distort the overall fairness and effectiveness of research assessment.
Advantages and disadvantages
Arguments in favor of open science generally focus on the value of increased transparency in research, and in the public ownership of science, particularly that which is publicly funded. In January 2014 J. Christopher Bare published a comprehensive "Guide to Open Science". Likewise, in 2017, a group of scholars known for advocating open science published a "manifesto" for open science in the journal Nature.
Advantages
Open access publication of research reports and data allows for rigorous peer-review
An article published by a team of NASA astrobiologists in 2010 in Science reported a bacterium known as GFAJ-1 that could purportedly metabolize arsenic (unlike any previously known species of lifeform). This finding, along with NASA's claim that the paper "will impact the search for evidence of extraterrestrial life", met with criticism within the scientific community. Much of the scientific commentary and critique around this issue took place in public forums, most notably on Twitter, where hundreds of scientists and non-scientists created a hashtag community around the hashtag #arseniclife. University of British Columbia astrobiologist Rosie Redfield, one of the most vocal critics of the NASA team's research, also submitted a draft of a research report of a study that she and colleagues conducted which contradicted the NASA team's findings; the draft report appeared in arXiv, an open-research repository, and Redfield called in her lab's research blog for peer review both of their research and of the NASA team's original paper. Researcher Jeff Rouder defined Open Science as "endeavoring to preserve the rights of others to reach independent conclusions about your data and work".
Publicly funded science will be publicly available
Public funding of research has long been cited as one of the primary reasons for providing Open Access to research articles. Since there is significant value in other parts of the research such as code, data, protocols, and research proposals a similar argument is made that since these are publicly funded, they should be publicly available under a Creative Commons Licence.
Open science will make science more reproducible and transparent
Increasingly the reproducibility of science is being questioned and for many papers or multiple fields of research was shown to be lacking. This problem has been described as a "reproducibility crisis". For example, psychologist Stuart Vyse notes that "(r)ecent research aimed at previously published psychology studies has demonstrated – shockingly – that a large number of classic phenomena cannot be reproduced, and the popularity of p-hacking is thought to be one of the culprits." Open Science approaches are proposed as one way to help increase the reproducibility of work as well as to help mitigate against manipulation of data.
Open science has more impact
There are several components to impact in research, many of which are hotly debated. However, under traditional scientific metrics parts Open science such as Open Access and Open Data have proved to outperform traditional versions.
Open Science can provide learning opportunities
Open science needs to acknowledge and accommodate the heterogeneity of science, it provides an opportunities for different communities to learn from other communities. For example preregistration in quantitative sciences can benefit qualitative researchers to reduce researcher degrees of freedom, whereas positionality statements have been used to contextual researcher and research environment in qualitative can be used in order to combat reproducibility crisis in quantitative research. In addition, journals should be open to publishing these behaviours, using a guide to ease journal editors into open science.
Open science will help answer uniquely complex questions
Recent arguments in favor of Open Science have maintained that Open Science is a necessary tool to begin answering immensely complex questions, such as the neural basis of consciousness, or pandemics such as the COVID-19 pandemic. The typical argument propagates the fact that these type of investigations are too complex to be carried out by any one individual, and therefore, they must rely on a network of open scientists to be accomplished. By default, the nature of these investigations also makes this "open science" as "big science". It is thought that open science could support innovation and societal benefits, supporting and reinforcing research activities by enabling digital resources that could, for example, use or provide structured open data.
Disadvantages
Arguments against open science tend to focus on the advantages of data ownership and concerns about the misuse of data, but see
Potential misuse
In 2011, Dutch researchers announced their intention to publish a research paper in the journal Science describing the creation of a strain of H5N1 influenza which can be easily passed between ferrets, the mammals which most closely mimic the human response to the flu. The announcement triggered a controversy in both political and scientific circles about the ethical implications of publishing scientific data which could be used to create biological weapons. These events are examples of how science data could potentially be misused. It has been argued that constraining the dissemination of dual-use knowledge can in certain cases be justified because, for example, "scientists have a responsibility for potentially harmful consequences of their research; the public need not always know of all scientific discoveries [or all its details]; uncertainty about the risks of harm may warrant precaution; and expected benefits do not always outweigh potential harm".
Scientists have collaboratively agreed to limit their own fields of inquiry on occasions such as the Asilomar conference on recombinant DNA in 1975, and a proposed 2015 worldwide moratorium on a human-genome-editing technique. Differential technological development aims to decrease risks by influencing the sequence in which technologies are developed. Relying only on the established form of legislation and incentives to ensure the right outcomes may not be adequate as these may often be too slow.
The public may misunderstand science data
In 2009 NASA launched the Kepler spacecraft and promised that they would release collected data in June 2010. Later they decided to postpone release so that their scientists could look at it first. Their rationale was that non-scientists might unintentionally misinterpret the data, and NASA scientists thought it would be preferable for them to be familiar with the data in advance so that they could report on it with their level of accuracy.
Low-quality science
Post-publication peer review, a staple of open science, has been criticized as promoting the production of lower quality papers that are extremely voluminous. Specifically, critics assert that as quality is not guaranteed by preprint servers, the veracity of papers will be difficult to assess by individual readers. This will lead to rippling effects of false science, akin to the recent epidemic of false news, propagated with ease on social media websites. Common solutions to this problem have been cited as adaptations of a new format in which everything is allowed to be published but a subsequent filter-curator model is imposed to ensure some basic quality of standards are met by all publications.
Entrapment by platform capitalism
For Philip Mirowski open science runs the risk of continuing a trend of commodification of science which ultimately serves the interests of capital in the guise of platform capitalism.
WEIRD-focus
Open Science is primarily driven by Western, Educated, Industrialized, Rich and Democratic (WEIRD) society that it is challenging for people from the Global South to implement or follow these changes for Open Science. As a result, it perpetuates inequalities found across cultures. However, journal editors have taken note of guidelines for change (e.g. ) in order to make sure Open Science is more inclusive with a focus of multi-site studies and value of diversity within Open Science discussion.
Actions and initiatives
Open-science projects
Different projects conduct, advocate, develop tools for, or fund open science.
The Allen Institute for Brain Science conducts numerous open science projects while the Center for Open Science has projects to conduct, advocate, and create tools for open science. Other workgroups have been created in different fields, such as the Decision Analysis in R for Technologies in Health (DARTH) workgroup], which is a multi-institutional, multi-university collaborative effort by researchers who have a common goal to develop transparent and open-source solutions to decision analysis in health.
Organizations have extremely diverse sizes and structures. The Open Knowledge Foundation (OKF) is a global organization sharing large data catalogs, running face to face conferences, and supporting open source software projects. In contrast, Blue Obelisk is an informal group of chemists and associated cheminformatics projects. The tableau of organizations is dynamic with some organizations becoming defunct, e.g., Science Commons, and new organizations trying to grow, e.g., the Self-Journal of Science. Common organizing forces include the knowledge domain, type of service provided, and even geography, e.g., OCSDNet's concentration on the developing world.
The Allen Brain Atlas maps gene expression in human and mouse brains; the Encyclopedia of Life documents all the terrestrial species; the Galaxy Zoo classifies galaxies; the International HapMap Project maps the haplotypes of the human genome; the Monarch Initiative makes available integrated public model organism and clinical data; and the Sloan Digital Sky Survey which regularizes and publishes data sets from many sources. All these projects accrete information provided by many different researchers with different standards of curation and contribution.
Mathematician Timothy Gowers launched open science journal Discrete Analysis in 2016 to demonstrate that a high-quality mathematics journal could be produced outside the traditional academic publishing industry. The launch followed a boycott of scientific journals that he initiated. The journal is published by a nonprofit which is owned and published by a team of scholars.
Other projects are organized around completion of projects that require extensive collaboration. For example, OpenWorm seeks to make a cellular level simulation of a roundworm, a multidisciplinary project. The Polymath Project seeks to solve difficult mathematical problems by enabling faster communications within the discipline of mathematics. The Collaborative Replications and Education project recruits undergraduate students as citizen scientists by offering funding. Each project defines its needs for contributors and collaboration.
Another practical example for open science project was the first "open" doctoral thesis started in 2012. It was made publicly available as a self-experiment right from the start to examine whether this dissemination is even possible during the productive stage of scientific studies. The goal of the dissertation project: Publish everything related to the doctoral study and research process as soon as possible, as comprehensive as possible and under an open license, online available at all time for everyone. End of 2017, the experiment was successfully completed and published in early 2018 as an open access book.
An example promoting accessibility of open-source code for research papers is CatalyzeX, which finds and links both official implementations by authors and source code independently replicated by other researchers. These code implementations are also surfaced on the preprint server arXiv and open peer-review platform OpenReview.
The ideas of open science have also been applied to recruitment with jobRxiv, a free and international job board that aims to mitigate imbalances in what different labs can afford to spend on hiring.
Advocacy
Numerous documents, organizations, and social movements advocate wider adoption of open science. Statements of principles include the Budapest Open Access Initiative from a December 2001 conference and the Panton Principles. New statements are constantly developed, such as the Amsterdam Call for Action on Open Science to be presented to the Dutch Presidency of the Council of the European Union in late May 2016. These statements often try to regularize licenses and disclosure for data and scientific literature.
Other advocates concentrate on educating scientists about appropriate open science software tools. Education is available as training seminars, e.g., the Software Carpentry project; as domain specific training materials, e.g., the Data Carpentry project; and as materials for teaching graduate classes, e.g., the Open Science Training Initiative. Many organizations also provide education in the general principles of open science.
Within scholarly societies there are also sections and interest groups that promote open science practices. The Ecological Society of America has an Open Science Section. Similarly, the Society for American Archaeology has an Open Science Interest Group.
Journal support
Many individual journals are experimenting with the open access model: the Public Library of Science, or PLOS, is creating a library of open access journals and scientific literature. Other publishing experiments include delayed and hybrid models. There are experiments in different fields:
F1000Research provides open publishing and open peer review for the life sciences.
The Open Library of Humanities is a non-profit open access publisher for the humanities and social sciences.
The Journals Library of the National Institute for Health and Care Research (NIHR) publishes all relevant documents and data from the onset of research projects, updating them alongside the progress of the study.
Journal support for open-science does not conflict with preprint servers:
figshare archives and shares images, readings, and other data; and Open Science Framework preprints, arXiv, and HAL Archives Ouvertes provide electronic preprints across many fields.
Software
A variety of computer resources support open science. These include software like the Open Science Framework from the Center for Open Science to manage project information, data archiving and team coordination; distributed computing services like Ibercivis to use unused CPU time for computationally intensive tasks; and services like Experiment.com to provide crowdsourced funding for research projects.
Blockchain platforms for open science have been proposed. The first such platform is the Open Science Organization, which aims to solve urgent problems with fragmentation of the scientific ecosystem and difficulties of producing validated, quality science. Among the initiatives of Open Science Organization include the Interplanetary Idea System (IPIS), Researcher Index (RR-index), Unique Researcher Identity (URI), and Research Network. The Interplanetary Idea System is a blockchain based system that tracks the evolution of scientific ideas over time. It serves to quantify ideas based on uniqueness and importance, thus allowing the scientific community to identify pain points with current scientific topics and preventing unnecessary re-invention of previously conducted science. The Researcher Index aims to establish a data-driven statistical metric for quantifying researcher impact. The Unique Researcher Identity is a blockchain technology based solution for creating a single unifying identity for each researcher, which is connected to the researcher's profile, research activities, and publications. The Research Network is a social networking platform for researchers. A scientific paper from November 2019 examined the suitability of blockchain technology to support open science.
Preprint servers
Preprint Servers come in many varieties, but the standard traits across them are stable: they seek to create a quick, free mode of communicating scientific knowledge to the public. Preprint servers act as a venue to quickly disseminate research and vary on their policies concerning when articles may be submitted relative to journal acceptance. Also typical of preprint servers is their lack of a peer-review process – typically, preprint servers have some type of quality check in place to ensure a minimum standard of publication, but this mechanism is not the same as a peer-review mechanism. Some preprint servers have explicitly partnered with the broader open science movement. Preprint servers can offer service similar to those of journals, and Google Scholar indexes many preprint servers and collects information about citations to preprints. The case for preprint servers is often made based on the slow pace of conventional publication formats. The motivation to start SocArXiv, an open-access preprint server for social science research, is the claim that valuable research being published in traditional venues often takes several months to years to get published, which slows down the process of science significantly. Another argument made in favor of preprint servers like SocArXiv is the quality and quickness of feedback offered to scientists on their pre-published work. The founders of SocArXiv claim that their platform allows researchers to gain easy feedback from their colleagues on the platform, thereby allowing scientists to develop their work into the highest possible quality before formal publication and circulation. The founders of SocArXiv further claim that their platform affords the authors the greatest level of flexibility in updating and editing their work to ensure that the latest version is available for rapid dissemination. The founders claim that this is not traditionally the case with formal journals, which instate formal procedures to make updates to published articles. Perhaps the strongest advantage of some preprint servers is their seamless compatibility with Open Science software such as the Open Science Framework. The founders of SocArXiv claim that their preprint server connects all aspects of the research life cycle in OSF with the article being published on the preprint server. According to the founders, this allows for greater transparency and minimal work on the authors' part.
One criticism of pre-print servers is their potential to foster a culture of plagiarism. For example, the popular physics preprint server ArXiv had to withdraw 22 papers when it came to light that they were plagiarized. In June 2002, a high-energy physicist in Japan was contacted by a man called Ramy Naboulsi, a non-institutionally affiliated mathematical physicist. Naboulsi requested Watanabe to upload his papers on ArXiv as he was not able to do so, because of his lack of an institutional affiliation. Later, the papers were realized to have been copied from the proceedings of a physics conference. Preprint servers are increasingly developing measures to circumvent this plagiarism problem. In developing nations like India and China, explicit measures are being taken to combat it. These measures usually involve creating some type of central repository for all available pre-prints, allowing the use of traditional plagiarism detecting algorithms to detect the fraud . Nonetheless, this is a pressing issue in the discussion of pre-print servers, and consequently for open science.
| Physical sciences | Science basics | Basics and measurement |
6279588 | https://en.wikipedia.org/wiki/Extranuclear%20inheritance | Extranuclear inheritance | Extranuclear inheritance or cytoplasmic inheritance is the transmission of genes that occur outside the nucleus. It is found in most eukaryotes and is commonly known to occur in cytoplasmic organelles such as mitochondria and chloroplasts or from cellular parasites like viruses or bacteria.
Organelles
Mitochondria are organelles which function to transform energy as a result of cellular respiration. Chloroplasts are organelles which function to produce sugars via photosynthesis in plants and algae. The genes located in mitochondria and chloroplasts are very important for proper cellular function. The mitochondrial DNA and other extranuclear types of DNA replicate independently of the DNA located in the nucleus, which is typically arranged in chromosomes that only replicate one time preceding cellular division. The extranuclear genomes of mitochondria and chloroplasts however replicate independently of cell division. They replicate in response to a cell's increasing energy needs which adjust during that cell's lifespan. Since they replicate independently, genomic recombination of these genomes is rarely found in offspring, contrary to nuclear genomes in which recombination is common.
Mitochondrial diseases are inherited from the mother, not from the father. Mitochondria with their mitochondrial DNA are already present in the egg cell before it gets fertilized by a sperm. In many cases of fertilization, the head of the sperm enters the egg cell; leaving its middle part, with its mitochondria, behind. The mitochondrial DNA of the sperm often remains outside the zygote and gets excluded from inheritance.
Parasites
Extranuclear transmission of viral genomes and symbiotic bacteria is also possible. An example of viral genome transmission is perinatal transmission. This occurs from mother to fetus during the perinatal period, which begins before birth and ends about 1 month after birth. During this time viral material may be passed from mother to child in the bloodstream or breastmilk. This is of particular concern with mothers carrying HIV or hepatitis C viruses. Symbiotic cytoplasmic bacteria are also inherited in organisms such as insects and protists.
Types
Three general types of extranuclear inheritance exist.
Vegetative segregation results from random replication and partitioning of cytoplasmic organelles. It occurs with chloroplasts and mitochondria during mitotic cell divisions and results in daughter cells that contain a random sample of the parent cell's organelles. An example of vegetative segregation is with mitochondria of asexually replicating yeast cells.
Uniparental inheritance occurs in extranuclear genes when only one parent contributes organellar DNA to the offspring. A classic example of uniparental gene transmission is the maternal inheritance of human mitochondria. The mother's mitochondria are transmitted to the offspring at fertilization via the egg. The father's mitochondrial genes are not transmitted to the offspring via the sperm. Very rare cases which require further investigation have been reported of paternal mitochondrial inheritance in humans, in which the father's mitochondrial genome is found in offspring. Chloroplast genes can also inherit uniparentally during sexual reproduction. They are historically thought to inherit maternally, but paternal inheritance in many species is increasingly being identified. The mechanisms of uniparental inheritance from species to species differ greatly and are quite complicated. For instance, chloroplasts have been found to exhibit maternal, paternal and biparental modes even within the same species. In tobacco (Nicotiana tabacum), the mode of chloroplast inheritance is affected by the temperature and the enzymatic activity of an exonuclease during male gametogenesis.
Biparental inheritance occurs in extranuclear genes when both parents contribute organellar DNA to the offspring. It may be less common than uniparental extranuclear inheritance, and usually occurs in a permissible species only a fraction of the time. An example of biparental mitochondrial inheritance is in the yeast Saccharomyces cerevisiae. When two haploid cells of opposite mating type fuse they can both contribute mitochondria to the resulting diploid offspring.
Mutant mitochondria
Poky is a mutant of the fungus Neurospora crassa that has extranuclear inheritance. Poky is characterized by slow growth, a defect in mitochondrial ribosome assembly and deficiencies in several cytochromes. The studies of poky mutants were among the first to establish an extranuclear mitochondrial basis for inheritance of a particular genotype. It was initially found, using genetic crosses, that poky is maternally inherited. Subsequently, the primary defect in the poky mutants was determined to be a deletion in the mitochondrial DNA sequence encoding the small subunit of mitochondrial ribosomal RNA.
| Biology and health sciences | Genetics | Biology |
6281041 | https://en.wikipedia.org/wiki/Cat%20flea | Cat flea | The cat flea (scientific name Ctenocephalides felis) is an extremely common parasitic insect whose principal host is the domestic cat, although a high proportion of the fleas found on dogs also belong to this species. This is despite the widespread existence of a separate and well-established "dog" flea, Ctenocephalides canis. Cat fleas originated in Africa but can now be found globally. As humans began domesticating cats, the prevalence of the cat flea increased and it spread throughout the world.
Of the cat fleas, Ctenocephalides felis felis is the most common, although other subspecies do exist, including C. felis strongylus, C. orientis, and C. damarensis. Over 90% of fleas found on both dogs and cats are Ctenocephalides felis felis.
Overview
The cat flea belongs to the insect order Siphonaptera which in its adult stage is an obligatory hematophage. Adults of both sexes range from 1–2 mm long and are usually a reddish-brown colour, although the abdomens of gravid females often swell with eggs causing them to appear banded in cream and dark brown. Like all fleas, the cat flea is compressed laterally allowing it to slip between the sometimes dense hairs of its host just above the top layer of the skin, resulting in an extremely thin insect that may be difficult to observe even if the host's coat is pure white. Cat fleas are wingless.
The cat flea affects both the cat and the dog worldwide. The cat flea can also maintain its life cycle on other carnivores and on omnivores, but these are only chosen when more acceptable hosts become unavailable. Adult cat fleas do not willingly leave their hosts, and inter-animal transfer of adult fleas is rare except in animals that share sleeping quarters. A flea which becomes separated from its host will often die within hours from starvation. It has been found that mortality differs between male and female cat fleas when separated from the host. It was found that within two days all male cat fleas were dead, while females became inactive after three days.
In addition to their role as pests in dogs and cats, cat fleas are responsible for a number of diseases. They can cause flea bite dermatitis and the transmission of dog tapeworm to name a few.
Life cycle
Cat fleas are holometabolous (undergo complete metamorphosis) insects and therefore go through four life cycle stages of egg, larva, pupa, and imago (adult). Adult fleas must feed on blood before they can become capable of reproduction.
Flea populations are distributed with about 50% eggs, 35% larvae, 10% pupae, and 5% adults. Cat fleas may live up to two years.
Eggs
An adult gravid female flea that has consumed a full blood meal will begin to produce between 20 and 30 microscopic (0.5 mm) non-adhesive white ovoid eggs per day, laying them individually and continually at a rate of about one per hour until she dies (under ideal conditions it might be possible for her to produce between 2,000 and 8,000 eggs in her lifetime, though most only manage to produce around 100 before being consumed by their host during grooming activity). The eggs are dispersed freely into the environment. Within two to seven weeks a certain proportion will then hatch into larvae. Hatching is at its highest when temperature is 27 °C and humidity is greater than 50%.
Given that eggs are non-adhesive, they do not stick to the host (70% are lost from the host in the first 8 hours).
Larvae
The larva of the cat flea has a grub-like appearance and is ~2 mm in length. The larvae are negatively phototaxic/phototropic, avoiding light and hiding in the substrate around them. The larvae require adequate ambient moisture and warmth, and will die at temperatures near freezing. Cat fleas prefer soil moisture content between 1-10%. While in this developmental stage the larvae will feed on a variety of organic substances, but the most important dietary item for them is the crumbs of dried blood that continually fall like snow out of the haircoat of the host after it has been excreted by the adult fleas as fecal material. Thus, the adult flea population continually feeds the larval population in the animal's environment. Adult feces is an important part of the larval diet. When reared in the lab, flea larvae provided with adult feces have a higher survival rate (67%) than those provided diets of dried bovine blood (39%) or meat flour (55%).
Pupal stage
Flea larvae metamorphose through four stages before spinning a cocoon and entering the pupal stage. The cocoon is adhesive, and quickly acquires a coat of camouflage from surrounding dirt and dust. Pupation depends heavily on temperature and moisture, and takes a week or more to complete, though a fully pupated adult can remain inside of its cocoon in a state of semi-dormancy (called the "pupal window") awaiting signs of the presence of a host.
Adult
Newly emerged fleas use variations in light and shadow along with increases in warmth and CO2 to detect the presence of a potential host, and will jump to a new host within seconds of emerging from the cocoon. The new flea begins feeding on host blood within minutes.
Effects on the hosts
A few fleas on adult dogs or cats cause little harm unless the host becomes allergic to substances in the flea's saliva. There are 15 substances that can cause allergy in flea saliva. The disease that results from allergy is called flea allergy dermatitis. Small animals with large infestations can lose enough bodily fluid to fleas feeding that dehydration may result. Cat fleas also may be responsible for disease transmission through humans, and have been suspected as transmission agents of plague. Severe flea infestations can result in anemia due to blood loss.
Disease transmission
Cat fleas can transmit other parasites and infections to dogs and cats and also to humans. The most prominent of these are Bartonella, murine typhus, and atopic dermatitis. The tapeworm Dipylidium caninum can be transmitted when an immature flea is swallowed by pets or humans. In addition, cat fleas have been found to carry Borrelia burgdorferi, the etiologic agent of Lyme disease, but their ability to transmit the disease is unclear. Finally, cat fleas are vectors for Rickettsia felis.
Fortunately, cat flee can not transmit Yersinia pestis.
Prevention and treatment of flea-borne disease
Since more than three-quarters of a flea's life is spent somewhere other than on the host animal, it is not adequate to treat only the host; it is important also to treat the host's environment. Thorough vacuuming, washing linens in hot water, and treating all hosts in the immediate environment (the entire household, for example) is essential for successful eradication. These steps should be performed on a regular basis as the flea life cycle is complex. Treatment should be implemented every five to ten days. Pet safe insecticides may also be an option in treating a pet with fleas, and soap is sufficient as an insecticide for adult fleas.
Insecticide resistance
Cat fleas have developed insecticide resistance to many of the common insecticides used to control them environmentally, including carbamates, organophosphates, and pyrethroids. Additionally, it has been found that larvae are more resistant to certain insecticides than adults. Targets of juvenile hormone may be successful to limit growth in the larval stages. When administering insecticides to pets for flea treatment, it is critically important to finish the full dose to limit the spread of resistance.
Impact of climate change on the cat flea
Cat fleas are generally tolerant to a wide range of environmental conditions. As the climate warms, however, it is predicted that the tropical haplotype will displace the temperate haplotype. Climate change often drives changes in species range. In Australia, it is predicted that warming temperatures will drive the cat flea distribution south.
| Biology and health sciences | Insects: General | Animals |
6283381 | https://en.wikipedia.org/wiki/Alcyonacea | Alcyonacea | Alcyonacea are an order of sessile colonial cnidarians that are found throughout the oceans of the world, especially in the deep sea, polar waters, tropics and subtropics. Whilst not in a strict taxonomic sense, Alcyonacea are commonly known as soft corals. The term "soft coral" generally applies to organisms in the two orders Pennatulacea and Alcyonacea with their polyps embedded within a fleshy mass of coenenchymal tissue. Consequently, the term "gorgonian coral" is commonly handed to multiple species in the order Alcyonacea that produce a mineralized skeletal axis (or axial-like layer) composed of calcite and the proteinaceous material gorgonin only and corresponds to only one of several families within the formally accepted taxon Gorgoniidae (Scleractinia). These can be found in order Malacalcyonacea (taxonomic synonyms of include (unnacepted): Alcyoniina, Holaxonia, Protoalcyonaria, Scleraxonia, and Stolonifera.
Common names for subsets of this order are sea fans and sea whips; others are similar to the sea pens of related order Pennatulacea. Individual tiny polyps form colonies that are normally erect, flattened, branching, and reminiscent of a fan. Others may be whiplike, bushy, or even encrusting. A colony can be several feet high and across, but only a few inches thick. They may be brightly coloured, often purple, red, or yellow. Photosynthetic gorgonians can be successfully kept in captive aquaria.
About 500 different species of gorgonians are found in the oceans of the world, but they are particularly abundant in the shallow waters of the Western Atlantic, including Florida, Bermuda, and the West Indies.
Anatomy
The structure of a gorgonian colony varies. In the suborder Holaxonia, skeletons are formed from a flexible, horny substance called gorgonin. The suborder Scleraxonia species are supported by a skeleton of tightly grouped calcareous spicules. Also, some species encrust like coral.
Measurements of the gorgonin and calcite within several long-lived species of gorgonians can be useful in paleoclimatology and paleoceanography, as their skeletal growth rate and composition are highly correlated with seasonal and climatic variation.
Features
Soft corals contain minute, spiny skeletal elements called sclerites, useful in species identification. Sclerites give these corals some degree of support and give their flesh a spiky, grainy texture that deters predators. In the past, soft corals were thought to be unable to lay new foundations for future corals, but recent findings suggest that colonies of the leather-coral genus Sinularia are able to cement sclerites and consolidate them at their base into alcyonarian spiculite, thus making them reef builders.
Unlike stony corals, most soft corals thrive in nutrient-rich waters with less intense light. Almost all use symbiotic photosynthetic zooxanthella as a major energy source. However, most readily eat any free-floating food, such as zooplankton, out of the water column. They are integral members of the reef ecosystem and provide habitat for fish, snails, algae, and a diversity of other marine species.
Despite being dominated by "soft corals", the order Alcyonacea now contains all species known as "gorgonian corals", that produce a hard skeleton made from gorgonin, a protein unique to the group that makes their skeletons quite different from "true" corals (Scleractinia). These "gorgonion corals" can be found in suborders Holaxonia, Scleraxonia, and Stolonifera.
Many soft corals are easily collected in the wild for the reef aquarium hobby, as small cuttings are less prone to infection or damage during shipping than stony corals. Nevertheless, home-grown specimens tend to be more adaptable to aquarium life and help conserve wild reefs. Soft corals grow quickly in captivity and are easily divided into new individuals, and so those grown by aquaculture are often hardier and less expensive than imported corals from the wild.
Ecology
Each gorgonian polyp has eight tentacles, which catch plankton and particulate matter for consumption. This process, called filter feeding, is facilitated when the "fan" is oriented across the prevailing current to maximise water flow to the gorgonian, hence food supply.
Some gorgonians contain algae, or zooxanthellae. This symbiotic relationship assists in giving the gorgonian nutrition by photosynthesis. Gorgonians possessing zooxanthellae are usually characterized by brownish polyps.
Gorgonians are found primarily in shallow waters, though some have been found at depths of several thousand feet. The size, shape, and appearance of gorgonians can be correlated with their location. The more fan-shaped and flexible gorgonians tend to populate shallower areas with strong currents, while the taller, thinner, and stiffer gorgonians can be found in deeper, calmer waters.
Other fauna, such as hydrozoa, bryozoa, and brittle stars, are known to dwell within the branches of gorgonian colonies. The pygmy seahorse not only makes certain species of gorgonians its home, but also closely resembles its host, thus is well camouflaged. Two species of pygmy seahorse, Hippocampus bargibanti and Hippocampus denise, are obligate residents on gorgonians. H. bargibanti is limited to two species in the single genus Muricella.
Gorgonians produce unusual organic compounds in their tissues, particularly diterpenes, and some of these are important candidates for new drugs. These compounds may be part of the chemical defenses produced by gorgonians to render their tissue distasteful to potential predators. Bottlenose dolphins in the Red Sea have been observed swimming against these tissues, in what is thought to be an attempt to take advantage of the antimicrobial qualities of diterpenes. Despite these chemical defenses, the tissues of gorgonians are prey for flamingo tongue snails of the genus Cyphoma, nudibranchs, the fireworm Hermodice spp., and their polyps are food for butterflyfishes. Amongst the nudibranchs which feed on soft corals and sea fans are the Tritoniidae and the genus Phyllodesmium which specialises in eating Xenia species.
Suborders and families
The World Register of Marine Species lists these suborders and families:
suborder Alcyoniina
family Acrophytidae McFadden & Ofwegen, 2017
family Alcyoniidae Lamouroux, 1812
family Aquaumbridae Breedy, van Ofwegen & Vargas, 2012
family Corymbophytidae McFadden & Ofwegen, 2017
family Leptophytidae McFadden & Ofwegen, 2017
family Nephtheidae Gray, 1862
family Nidaliidae Gray, 1869
family Paralcyoniidae Gray, 1869
family Xeniidae Ehrenberg, 1828
suborder Calcaxonia
family Chrysogorgiidae Verrill, 1883
family Ellisellidae Gray, 1859
family Ifalukellidae Bayer, 1955
family Isididae Lamouroux, 1812
family Primnoidae Milne Edwards, 1857
suborder Holaxonia
family Acanthogorgiidae Gray, 1859
family Dendrobrachiidae Brook, 1889
family Gorgoniidae Lamouroux, 1812
family Keroeididae Kinoshita, 1910
family Plexauridae Gray, 1859
suborder Protoalcyonaria
family Taiaroidae Bayer & Muzik, 1976
suborder Scleraxonia
family Anthothelidae Broch, 1916
family Briareidae Gray, 1859
family Coralliidae Lamouroux, 1812
family Melithaeidae Gray, 1870
family Paragorgiidae Kükenthal, 1916
family Parisididae Aurivillius, 1931
family Spongiodermidae Wright & Studer, 1889
family Subergorgiidae Gray, 1859
family Victorgorgiidae Moore, Alderslade & Miller, 2017
suborder Stolonifera
family Acrossotidae Bourne, 1914
family Arulidae McFadden & van Ofwegen, 2012
family Clavulariidae Hickson, 1894
family Coelogorgiidae Bourne, 1900
family Cornulariidae Dana, 1846
family Pseudogorgiidae Utinomi & Harada, 1973
family Tubiporidae Ehrenberg, 1828
family Acanthoaxiidae van Ofwegen & McFadden, 2010
family Haimeidae Wright, 1865
family Paramuriceidae Bayer, 1956
family Parasphaerascleridae McFadden & van Ofwegen, 2013
family Viguieriotidae
| Biology and health sciences | Cnidarians | Animals |
6286608 | https://en.wikipedia.org/wiki/Sequence%20%28geology%29 | Sequence (geology) | In geology, a sequence is a stratigraphic unit which is bounded by an unconformity at the top and at the bottom.
Definition
In a more rigorous and general way, a sequence is defined as a
"relatively conformable [...], genetically related succession of strata bounded by unconformities or their correlative surfaces"
Special cases and related concepts
Special cases of sequences include type 1 sequences and type 2 sequences. A related concept are parasequences. Contrary to their name they are not smaller sequences.
| Physical sciences | Stratigraphy | Earth science |
477823 | https://en.wikipedia.org/wiki/Gasterosteoidei | Gasterosteoidei | Gasterosteoidei is a suborder of ray-finned fishes that includes the sticklebacks and relatives, the 5th edition of Fishes of the World classifies this suborder within the order Scorpaeniformes.
Systematics
Gasterosteoidei is treated as a suborder within the order Scorpaeniformes in the 5th edition of Fishes of the World, but in other phylogenetic classifications it is treated as the infraorder Gasterosteales within the suborder Cottoidei or as a sister clade to the Zoarcales in the order Zoarciformes. Indostomidae is included within Gasterosteoidei in Fishes of the World''' but according to Betancur et al its inclusion in the clade renders it paraphyletic and they classify that family within the monotypic suborder Indostomoidei within the Synbranchiformes.
Historically, Gasterosteoidei was treated as a suborder within the order Gasterostiformes and often included the sea horses, pipefishes and their relatives as suborder Syngnathoidei, with the sticklebacks and relatives in the suborder Gasterosteoidei. The Gasterosteiformes sensu lato were regarded as paraphyletic with the Scorpaeniformes. The more typical members of that group (e.g. scorpionfishes) are apparently closer to the "true" Gasterosteiformes, whereas the keel-bodied flying gurnards (Dactylopteridae) seem actually to belong to the Syngnathiformes clade. It seems that the closest living relatives of the narrowly delimited Gasterosteoidei are the Zoarcoidei, which have been placed in the massively paraphyletic "Perciformes". The Zoarcoidei, as well as the related Trichodontidae, would then appear to be derived offshoots of the scorpaeniform-gasterosteiform radiation which have apomorphically lost the bone "armour" found in their relatives.
Families and genera
Gasterosteoidei contains the following families and genera:
Family Hypoptychidae Steindachner, 1880 (Sand eel)
Hypoptychus Steindachner, 1880
Family Aulorhynchidae Gill, 1861 (Tubesnouts)
Aulichthys Brevoort 1862
Aulorhynchus Gill, 1861
Family Gasterosteidae Bonaparte, 1831 (Sticklebacks)
Apeltes DeKay, 1842
Culaea Whitley, 1950
Gasterosteus Linnaeus, 1758
Pungitius d'Annone, 1760
Spinachia Cuvier, 1816
Characteristics
Gasterosteoidei is characterised by the possession of a protractile upper jaw and a well developed upward pointing process on the premaxilla. The body is often armoured with dermal plates and paired dermal plates grow from membranes growing out fronm the pelvic girdle. If there are plates on the flanks these are often a single row of ossified lateral and dermal plates. Unpaired plates
paired pelvic plates arising from a membranous outgrowth of the pelvic girdle; lateral body plates, when present, are represented by a single series of lateral and dermal ossifications. The unpaired plates on the body which create the dorsal and ventral series grow from the expanded proximal middle radials of the pterygiphores of the dorsal and anal fins. Separate pectoral radials do not develop during the fish's development and the pectoral radial plate is fused into a single unit on the scapulo-coracoid. They have very small mouths. There are between 1 and 6 branchiostegal rays and there is no postcleithrum in the pelvic girdle which is never joined directly to the cleithra. There are other skeletal features that these fishes share too. The kidneys of gasterosteoids synthesis an adhesive chemical which is used by males to create nests of plant material, it is not known if this is true of all the taxa within the group. These are all rather small fishes with the largest species being the sea stickleback (Spinachia spinachia'') which has a maximum published standard length of .
Distribution and habitat
Gasterodteoidei are found in the northern hemisphere, mostly within the temperate and Arctic regions, the exception is the Indostomidae which are found in freshwater habitats in mainland Southeast Asia. The other groups can be found in fresh, brackish and salt water.
Timeline of genera
Source:
| Biology and health sciences | Acanthomorpha | Animals |
477903 | https://en.wikipedia.org/wiki/Foraminifera | Foraminifera | Foraminifera ( ; Latin for "hole bearers"; informally called "forams") are single-celled organisms, members of a phylum or class of Rhizarian protists characterized by streaming granular ectoplasm for catching food and other uses; and commonly an external shell (called a "test") of diverse forms and materials. Tests of chitin (found in some simple genera, and Textularia in particular) are believed to be the most primitive type. Most foraminifera are marine, the majority of which live on or within the seafloor sediment (i.e., are benthic, with different sized species playing a role within the macrobenthos, meiobenthos, and microbenthos), while a smaller number float in the water column at various depths (i.e., are planktonic), which belong to the suborder Globigerinina. Fewer are known from freshwater or brackish conditions, and some very few (nonaquatic) soil species have been identified through molecular analysis of small subunit ribosomal DNA.
Foraminifera typically produce a test, or shell, which can have either one or multiple chambers, some becoming quite elaborate in structure. These shells are commonly made of calcium carbonate () or agglutinated sediment particles. Over 50,000 species are recognized, both living (6,700–10,000) and fossil (40,000). They are usually less than 1 mm in size, but some are much larger, the largest species reaching up to 20 cm.
In modern scientific English, the term foraminifera is both singular and plural (irrespective of the word's Latin derivation), and is used to describe one or more specimens or taxa: its usage as singular or plural must be determined from context. Foraminifera is frequently used informally to describe the group, and in these cases is generally lowercase.
History of study
The earliest known reference to foraminifera comes from Herodotus, who in the 5th century BCE noted them as making up the rock that forms the Great Pyramid of Giza. These are today recognized as representatives of the genus Nummulites. Strabo, in the 1st Century BCE, noted the same foraminifera, and suggested that they were the remains of lentils left by the workers who built the pyramids.
Robert Hooke observed a foraminifera under the microscope, as described and illustrated in his 1665 book Micrographia:I was trying several small and single Magnifying Glasses, and casually viewing a parcel of white Sand, when I perceiv'd one of the grains exactly shap'd and wreath'd like a Shell[...] I view'd it every way with a better Microscope and found it on both sides, and edge-ways, to resemble the Shell of a small Water-Snail with a flat spiral Shell[...]Antonie van Leeuwenhoek described and illustrated foraminiferal tests in 1700, describing them as minute cockles; his illustration is recognizable as being Elphidium. Early workers classified foraminifera within the genus Nautilus, noting their similarity to certain cephalopods. It was recognised by Lorenz Spengler in 1781 that foraminifera had holes in the septa, which would eventually grant the group its name. Spengler also noted that the septa of foraminifera arced the opposite way from those of nautili and that they lacked a nerve tube.
Alcide d'Orbigny, in his 1826 work, considered them to be a group of minute cephalopods and noted their odd morphology, interpreting the pseudopodia as tentacles and noting the highly reduced (in actuality, absent) head. He named the group foraminifères, or "hole-bearers", as members of the group had holes in the divisions between compartments in their shells, in contrast to nautili or ammonites.
The protozoan nature of foraminifera was first recognized by Dujardin in 1835. Shortly after, in 1852, d'Orbigny produced a classification scheme, recognising 72 genera of foraminifera, which he classified based on test shape—a scheme that drew severe criticism from colleagues.
H.B. Brady's 1884 monograph described the foraminiferal finds of the Challenger expedition. Brady recognized 10 families with 29 subfamilies, with little regard to stratigraphic range; his taxonomy emphasized the idea that multiple different characters must separate taxonomic groups, and as such placed agglutinated and calcareous genera in close relation.
This overall scheme of classification would remain until Cushman's work in the late 1920s. Cushman viewed wall composition as the single most important trait in classification of foraminifera; his classification became widely accepted but also drew criticism from colleagues for being "not biologically sound".
Geologist Irene Crespin undertook extensive research in this field, publishing some ninety papers—including notable work on foraminifera—as sole author as well as more than twenty in collaboration with other scientists.
Cushman's scheme nevertheless remained the dominant scheme of classification until Tappan and Loeblich's 1964 classification, which placed foraminifera into the general groupings still used today, based on microstructure of the test wall. These groups have been variously moved around according to different schemes of higher-level classification. Pawlowski's (2013) use of molecular systematics has generally confirmed Tappan and Loeblich's groupings, with some being found as polyphyletic or paraphyletic; this work has also helped to identify higher-level relationships among major foraminiferal groups.
Taxonomy
The taxonomic position of the Foraminifera has varied since Schultze in 1854,
who referred to as an order, Foraminiferida. Loeblich (1987) and Tappan (1992) reranked Foraminifera as a class as it is now commonly regarded.
The Foraminifera have typically been included in the Protozoa, or in the similar Protoctista or Protist kingdom. Compelling evidence, based primarily on molecular phylogenetics, exists for their belonging to a major group within the Protozoa known as the Rhizaria. Prior to the recognition of evolutionary relationships among the members of the Rhizaria, the Foraminifera were generally grouped with other amoeboids as phylum Rhizopodea (or Sarcodina) in the class Granuloreticulosa.
The Rhizaria are problematic, as they are often called a "supergroup", rather than using an established taxonomic rank such as phylum. Cavalier-Smith defines the Rhizaria as an infra-kingdom within the kingdom Protozoa.
Some taxonomies put the Foraminifera in a phylum of their own, putting them on par with the amoeboid Sarcodina in which they had been placed.
Although as yet unsupported by morphological correlates, molecular data strongly suggest the Foraminifera are closely related to the Cercozoa and Radiolaria, both of which also include amoeboids with complex shells; these three groups make up the Rhizaria. However, the exact relationships of the forams to the other groups and to one another are still not entirely clear. Foraminifera are closely related to testate amoebae.
Anatomy
The most striking aspect of most foraminifera are their hard shells, or tests. These may consist of one of multiple chambers, and may be composed of protein, sediment particles, calcite, aragonite, or (in one case) silica. Some foraminifera lack tests entirely. Unlike other shell-secreting organisms, such as molluscs or corals, the tests of foraminifera are located inside the cell membrane, within the protoplasm. The organelles of the cell are located within the of the test, and the of the test allow the transfer of material from the pseudopodia to the internal cell and back.
The foraminiferal cell is divided into granular endoplasm and transparent ectoplasm from which a pseudopodial net may emerge through a single opening or through many perforations in the test. Individual pseudopods characteristically have small granules streaming in both directions. Foraminifera are unique in having granuloreticulose pseudopodia; that is, their pseudopodia appear granular under the microscope; these pseudopodia are often elongate and may split and rejoin each other. These can be extended and retracted to suit the needs of the cell. The pseudopods are used for locomotion, anchoring, excretion, test construction and in capturing food, which consists of small organisms such as diatoms or bacteria.
Aside from the tests, foraminiferal cells are supported by a cytoskeleton of microtubules, which are loosely arranged without the structure seen in other amoeboids. Forams have evolved special cellular mechanisms to quickly assemble and disassemble microtubules, allowing for the rapid formation and retraction of elongated pseudopodia.
In the gamont (sexual form), foraminifera generally have only a single nucleus, while the agamont (asexual form) tends to have multiple nuclei. In at least some species the nuclei are dimorphic, with the somatic nuclei containing three times as much protein and RNA than the generative nuclei. However, nuclear anatomy seems to be highly diverse. The nuclei are not necessarily confined to one chamber in multi-chambered species. Nuclei can be spherical or have many lobes. Nuclei are typically 30-50 μm in diameter.
Some species of foraminifera have large, empty vacuoles within their cells; the exact purpose of these is unclear, but they have been suggested to function as a reservoir of nitrate.
Mitochondria are distributed evenly throughout the cell, though in some species they are concentrated under the pores and around the external margin of the cell. This has been hypothesised to be an adaptation to low-oxygen environments.
Several species of xenophyophore have been found to have unusually high concentrations of radioactive isotopes within their cells, among the highest of any eukaryote. The purpose of this is unknown.
Ecology
Modern Foraminifera are primarily marine organisms, but living individuals have been found in brackish, freshwater and even terrestrial habitats. The majority of the species are benthic, and a further 50 morphospecies are planktonic. This count may, however, represent only a fraction of actual diversity, since many genetically distinct species may be morphologically indistinguishable.
Benthic foraminifera are typically found in fine-grained sediments, where they actively move between layers; however, many species are found on hard rock substrates, attached to seaweeds, or sitting atop the sediment surface.
The majority of planktonic foraminifera are found in the globigerinina, a lineage within the rotaliida. However, at least one other extant rotaliid lineage, Neogallitellia, seems to have independently evolved a planktonic lifestyle. Further, it has been suggested that some Jurassic fossil foraminifera may have also independently evolved a planktonic lifestyle, and may be members of Robertinida.
A number of forams, both benthic and planktonic, have unicellular algae as endosymbionts, from diverse lineages such as the green algae, red algae, golden algae, diatoms, and dinoflagellates. These mixotrophic foraminifers are particularly common in nutrient-poor oceanic waters. Some forams are kleptoplastic, retaining chloroplasts from ingested algae to conduct photosynthesis.
Most foraminifera are heterotrophic, consuming smaller organisms and organic matter; some smaller species are specialised feeders on phytodetritus, while others specialise in consuming diatoms. Some benthic forams construct feeding cysts, using the pseuodopodia to encyst themselves inside of sediment and organic particles. Certain foraminifera prey upon small animals such as copepods or cumaceans; some forams even predate upon other forams, drilling holes into the tests of their prey. One group, the xenophyophores, has been suggested to farm bacteria within their tests, although studies have failed to find support for this hypothesis. Suspension feeding is also common in the group, and at least some species can take advantage of dissolved organic carbon.
A few foram species are parasitic, infecting sponges, molluscs, corals, or even other foraminifera. Parasitic strategies vary; some act as ectoparasites, using their pseudopodia to steal food from the host, while others burrow through the shell or body wall of their host to feed on its soft tissue.
Foraminifera are themselves eaten by a host of larger organisms, including invertebrates, fish, shorebirds, and other foraminifera. It has been suggested, however, that in some cases predators may be more interested in the calcium from foram shells than in the organisms themselves. Several aquatic snail species are known to selectively feed upon foraminifera, often even preferring individual species.
Certain benthic foraminifera have been found to be capable of surviving anoxic conditions for over 24 hours, indicating that they are capable of selective anaerobic respiration. This is interpreted as an adaptation to survive changing oxygenic conditions near the sediment-water interface.
Foraminifera are found in the deepest parts of the ocean such as the Mariana Trench, including the Challenger Deep, the deepest part known. At these depths, below the carbonate compensation depth, the calcium carbonate of the tests is soluble in water due to the extreme pressure. The Foraminifera found in the Challenger Deep thus have no carbonate test, but instead have one of organic material.
Nonmarine foraminifera have traditionally been neglected in foram research, but recent studies show them to be substantially more diverse than previously known. They are known to inhabit disparate ecological niches, including mosses, rivers, lakes and ponds, wetlands, soils, peat bogs, and sand dunes.
Reproduction
The generalized foraminiferal life-cycle involves an alternation between haploid and diploid generations, although they are mostly similar in form. The haploid or gamont initially has a single nucleus, and divides to produce numerous gametes, which typically have two flagella. The diploid or agamont is multinucleate, and after meiosis divides to produce new gamonts. Multiple rounds of asexual reproduction between sexual generations are not uncommon in benthic forms.
Foraminifera exhibit morphological dimorphism associated with their reproductive cycle. The gamont, or sexually reproducing haploid form, is megalospheric—that is, its proloculus, or first chamber, is proportionally large. The gamont is also known as the A form. Gamonts, despite having typically larger proloculi, also generally have smaller overall test diameter than do agamonts.
After reaching maturity, the gamont divides via mitosis to produce thousands of gametes which are also haploid. These gametes all have a full set of organelles, and are expelled from the test into the environment leaving the test undamaged. Gametes are not differentiated into sperm and egg, and any two gametes from a species can generally fertilize each other.
When two gametes combine, they create a diploid, multi-nucleated cell known as the agamont, or B form. In contrast to the gamont, the agamont is microspheric, with a proportionally small first chamber but typically larger overall diameter with more chambers. The agamont is the asexual reproduction phase of the foraminifera; upon reaching adulthood, the protoplasm entirely vacates the test and divides its cytoplasm meiotically via multiple fission to form a number of haploid offspring. These offspring then begin to form their megalospheric first chamber before dispersing.
In some cases the haploid young may mature into a megalospheric form which then reproduces asexually to produce another megalospheric, haploid offspring. In this case, the first megalospheric form is referred to as the schizont or A1 form, while the second is referred to as the gamont or A2 form.
Maturation and reproduction occur more slowly in cooler and deeper water; these conditions also cause forams to grow larger. A forms always seem to be much more numerous than are B forms, likely due to the reduced likelihood of two gametes encountering one another and successfully combining.
Variations in reproductive mode
There is a high degree of diversity in reproductive strategies in different foraminiferal groups.
In unilocular species, the A form and B form are still present. As in the microspheric morph of multilocular forams, the asexually reproducing B form is larger than the sexually reproducing A form.
Forams in the family Spirillinidae have amoeboid gametes rather than flagellated. Other aspects of reproduction in this group are generally similar to that of other groups of forams.
The calcareous spirillinid Patellina corrugata has a slightly different reproductive strategy than most other foraminifera. The asexually reproducing B form produces a cyst that surrounds the entire cell; it then divides within this cyst and the juvenile cells cannibalise the calcite of the parent's test to form the first chamber of their own test. These A forms, upon maturity, gather into groups of up to nine individuals; they then form a protective cyst around the whole group. Gametogenesis occurs within this cyst, producing very low numbers of gametes. The B form larvae are produced inside of the cyst; any nuclei that are not bound into cells are consumed as food for the developing larvae. Patellina in A form is reportedly dioecious, with sexes referred to as the "plus" and "minus"; these sexes differ in number of nuclei, with the "plus" form having three nuclei and the "minus" form having four nuclei. The B form is again larger than the A form.
Tests
Foraminiferal tests serve to protect the organism within. Owing to their generally hard and durable construction (compared to other protists), the tests of foraminifera are a major source of scientific knowledge about the group.
Openings in the test that allow the cytoplasm to extend outside are called apertures. The primary aperture, leading to the exterior, take many different shapes in different species, including but not limited to rounded, crescent-shaped, slit-shaped, hooded, radiate (star-shaped), dendritic (branching). Some foraminifera have "toothed", flanged, or lipped primary apertures. There may be only one primary aperture or multiple; when multiple are present, they may be clustered or equatorial. In addition to the primary aperture, many foraminifera have supplemental apertures. These may form as relict apertures (past primary apertures from an earlier growth stage) or as unique structures.
Test shape is highly variable among different foraminifera; they may be single-chambered (unilocular) or multi-chambered (multilocular). In multilocular forms, new chambers are added as the organism grows. A wide variety of test morphologies is found in both unilocular and multilocular forms, including spiraled, serial, and milioline, among others.
Many foraminifera exhibit dimorphism in their tests, with megalospheric and microspheric individuals. These names should not be taken as referring to the size of the full organism; rather, they refer to the size of the first chamber, or proloculus. Tests as fossils are known from as far back as the Ediacaran period, and many marine sediments are composed primarily of them. For instance, the limestone that makes up the pyramids of Egypt is composed almost entirely of nummulitic benthic Foraminifera. It is estimated that reef Foraminifera generate about 43 million tons of calcium carbonate per year.
Genetic studies have identified the naked amoeba Reticulomyxa and the peculiar xenophyophores as foraminiferans without tests. A few other amoeboids produce reticulose pseudopods, and were formerly classified with the forams as the Granuloreticulosa, but this is no longer considered a natural group, and most are now placed among the Cercozoa.
Evolutionary history
Molecular clocks indicate that the crown-group of foraminifera likely evolved during the Neoproterozoic, between 900 and 650 million years ago; this timing is consistent with Neoproterozoic fossils of the closely related filose amoebae. As fossils of foraminifera have not been found prior to the very end of the Ediacaran, it is likely that most of these Proterozoic forms did not have hard-shelled tests.
Due to their non-mineralised tests, "allogromiids" have no fossil record.
The mysterious vendozoans of the Ediacaran period have been suggested to represent fossil xenophyophores. However, the discovery of diagenetically altered C27 sterols associated with the remains of Dickinsonia cast doubt on this identification and suggest it may instead be an animal. Other researchers have suggested that the elusive trace fossil Paleodictyon and its relatives may represent a fossil xenophyophore and noted the similarity of the extant xenophyophore Occultammina to the fossil; however, modern examples of Paleodictyon have not been able to clear up the issue and the trace may alternately represent a burrow or a glass sponge. Supporting this notion is the similar habitat of living xenophyophores to the inferred habitat of fossil graphoglyptids; however, the large size and regularity of many graphoglyptids as well as the apparent absence of xenophyae in their fossils casts doubt on the possibility. As of 2017 no definite xenophyophore fossils have been found.
Test-bearing foraminifera have an excellent fossil record throughout the Phanerozoic eon. The earliest known definite foraminifera appear in the fossil record towards the very end of the Ediacaran; these forms all have agglutinated tests and are unilocular. These include forms like Platysolenites and Spirosolenites.
Single-chambered foraminifera continued to diversify throughout the Cambrian. Some commonly encountered forms include Ammodiscus, Glomospira, Psammosphera, and Turritellella; these species are all agglutinated. They make up part of the Ammodiscina, a lineage of spirillinids that still contains modern forms. Later spirillinids would evolve multilocularity and calcitic tests, with the first such forms appearing during the Triassic; the group saw little effects on diversity due to the K-Pg extinction.
The earliest multi-chambered foraminifera are agglutinated species, and appear in the fossil record during the middle Cambrian period. Due to their poor preservation they cannot be positively assigned to any major foram group.
The earliest known calcareous-walled foraminifera are the Fusulinids, which appear in the fossil record during the Llandoverian epoch of the early Silurian. The earliest of these were microscopic, planispirally coiled, and evolute; later forms evolved a diversity of shapes including lenticular, globular, and elongated rice-shaped forms.
Later species of fusulinids grew to much larger size, with some forms reaching 5 cm in length; reportedly, some specimens reach up to 14 cm in length, making them among the largest foraminifera extant or extinct. Fusulinids are the earliest lineage of foraminifera thought to have evolved symbiosis with photosynthetic organisms. Fossils of fusulinids have been found on all continents except Antarctica; they reached their greatest diversity during the Visean epoch of the Carboniferous. The group then gradually declined in diversity until finally going extinct during the Permo-Triassic extinction event.
During the Tournaisian epoch of the Carboniferous, Miliolid foraminifera first appeared in the fossil record, having diverged from the spirillinids within the Tubothalamea. Miliolids suffered about 50% casualties during both the Permo-Triassic and K-Pg extinctions but survived to the present day. Some fossil miliolids reached up to 2 cm in diameter.
The earliest known Lagenid fossils appear during the Moscovian epoch of the Carboniferous. Seeing little effect due to the Permo-Triassic or K-Pg extinctions, the group diversified through time. Secondarily unilocular taxa evolved during the Jurassic and Cretaceous.
The earliest Involutinid fossils appear during the Permian; the lineage diversified throughout the Mesozoic of Eurasia before apparently vanishing from the fossil record following the Cenomanian-Turonian Ocean Anoxic Event. The extant group planispirillinidae has been referred to the involutinida, but this remains the subject of debate.
The Robertinida first appear in the fossil record during the Anisian epoch of the Triassic. The group remained at low diversity throughout its fossil history; all living representatives belong to the Robertinidae, which first appeared during the Paleocene.
The first definite Rotaliid fossils do not appear in the fossil record until the Pliensbachian epoch of the Jurassic, following the Triassic-Jurassic event. Diversity of the group remained low until the aftermath of the Cenomanian-Turonian event, after which the group saw a rapid diversification. Of this group, the planktonic Globigerinina—the first known group of planktonic forams—first appears in the aftermath of the Toarcian Turnover; the group saw heavy losses during both the K-Pg extinction and the Eocene-Oligocene extinction, but remains extant and diverse to this day. An additional evolution of planktonic lifestyle occurred in the Miocene or Pliocene, when the rotaliid Neogallitellia independently evolved a planktonic lifestyle.
Paleontological applications
Dying planktonic Foraminifera continuously rain down on the sea floor in vast numbers, their mineralized tests preserved as fossils in the accumulating sediment. Beginning in the 1960s, and largely under the auspices of the Deep Sea Drilling, Ocean Drilling, and International Ocean Drilling Programmes, as well as for the purposes of oil exploration, advanced deep-sea drilling techniques have been bringing up sediment cores bearing Foraminifera fossils. The effectively unlimited supply of these fossil tests and the relatively high-precision age-control models available for cores has produced an exceptionally high-quality planktonic Foraminifera fossil record dating back to the mid-Jurassic, and presents an unparalleled record for scientists testing and documenting the evolutionary process. The exceptional quality of the fossil record has allowed an impressively detailed picture of species inter-relationships to be developed on the basis of fossils, in many cases subsequently validated independently through molecular genetic studies on extant specimens
Because certain types of foraminifera are found only in certain environments, their fossils can be used to figure out the kind of environment under which ancient marine sediments were deposited; conditions such as salinity, depth, oxygenic conditions, and light conditions can be determined from the different habitat preferences of various species of forams. This allows workers to track changing climates and environmental conditions over time by aggregating information about the foraminifera present.
In other cases, the relative proportion of planktonic to benthic foraminifera fossils found in a rock can be used as a proxy for the depth of a given locality when the rocks were being deposited.
Since at least 1997, the Paleocene–Eocene thermal maximum (PETM) has been investigated as an analogy for understanding the effects of global warming and of massive carbon inputs to the ocean and atmosphere, including ocean acidification. Humans today emit about 10 Gt of carbon (about 37 Gt CO2e) per year, and at that rate will release a comparable amount to the PETM in about one thousand years. A main difference is that during the PETM the planet was ice-free, as the Drake Passage had not yet opened and the Central American Seaway had not yet closed. Although the PETM is now commonly held to be a case study for global warming and massive carbon emission, the cause, details, and overall significance of the event remain uncertain.
Foraminifera have significant application in the field of biostratigraphy. Due to their small size and hard shells, foraminifera may be preserved in great abundance and with high quality of preservation; due to their complex morphology, individual species are easily recognizable. Foraminifera species in the fossil record have limited ranges between the species' first evolution and their disappearance; stratigraphers have worked out the successive changes in foram assemblages throughout much of the Phanerozoic. As such, the assemblage of foraminifera within a given locality can be analyzed and compared to known dates of appearance and disappearance in order to narrow down the age of the rocks. This allows paleontologists to interpret the age of sedimentary rocks when radiometric dating is not applicable. This application of foraminifera was discovered by Alva C. Ellisor in 1920.
Calcareous fossil foraminifera are formed from elements found in the ancient seas where they lived. Thus, they are very useful in paleoclimatology and paleoceanography. They can be used, as a climate proxy, to reconstruct past climate by examining the stable isotope ratios and trace element content of the shells (tests). Global temperature and ice volume can be revealed by the isotopes of oxygen, and the history of the carbon cycle and oceanic productivity by examining the stable isotope ratios of carbon; see δ18O and δ13C. The concentration of trace elements, like strontium (Sr), magnesium (Mg), lithium (Li) and boron (B), also hold a wealth of information about global temperature cycles, continental weathering, and the role of the ocean in the global carbon cycle. Geographic patterns seen in the fossil records of planktonic forams are also used to reconstruct ancient ocean currents.
Modern uses
The oil industry relies heavily on microfossils such as forams to find potential hydrocarbon deposits. For the same reasons they make useful biostratigraphic markers, living foraminiferal assemblages have been used as bioindicators in coastal environments, including indicators of coral reef health. Because calcium carbonate is susceptible to dissolution in acidic conditions, foraminifera may be particularly affected by changing climate and ocean acidification.
Foraminifera have many uses in petroleum exploration and are used routinely to interpret the ages and paleoenvironments of sedimentary strata in oil wells. Agglutinated fossil foraminifera buried deeply in sedimentary basins can be used to estimate thermal maturity, which is a key factor for petroleum generation. The Foraminiferal Colouration Index (FCI) is used to quantify colour changes and estimate burial temperature. FCI data is particularly useful in the early stages of petroleum generation (about 100 °C).
Foraminifera can also be used in archaeology in the provenancing of some stone raw material types. Some stone types, such as limestone, are commonly found to contain fossilised foraminifera. The types and concentrations of these fossils within a sample of stone can be used to match that sample to a source known to contain the same "fossil signature".
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477917 | https://en.wikipedia.org/wiki/Lemur | Lemur | Lemurs ( ; from Latin ) are wet-nosed primates of the superfamily Lemuroidea ( ), divided into 8 families and consisting of 15 genera and around 100 existing species. They are endemic to the island of Madagascar. Most existing lemurs are small, have a pointed snout, large eyes, and a long tail. They chiefly live in trees and are active at night.
Lemurs share resemblance with other primates, but evolved independently from monkeys and apes. Due to Madagascar's highly seasonal climate, lemur evolution has produced a level of species diversity rivaling that of any other primate group.
Living lemurs range in weight from the mouse lemur to the indri. Since the arrival of humans on the island around 2,000 years ago, over a dozen species of "giant lemurs" larger than living lemur species have become extinct, including the gorilla-sized Archaeoindris. Lemurs share many common basal primate traits, such as divergent digits on their hands and feet, and nails instead of claws (in most species). However, their brain-to-body size ratio is smaller than that of anthropoid primates. As with all strepsirrhine primates, they have a "wet nose" (rhinarium).
Lemurs are generally the most social of the strepsirrhine primates, living in groups known as troops. They communicate more with scents and vocalizations than with visual signals. Lemurs have a relatively low basal metabolic rate, and as a result may exhibit dormancy such as hibernation or torpor. They also have seasonal breeding and female social dominance. Most eat a wide variety of fruits and leaves, while some are specialists. Two species of lemurs may coexist in the same forest due to different diets.
Lemur research during the 18th and 19th centuries focused on taxonomy and specimen collection. Modern studies of lemur ecology and behavior did not begin in earnest until the 1950s and 1960s. Initially hindered by political issues on Madagascar during the mid-1970s, field studies resumed in the 1980s. Lemurs are important for research because their mix of ancestral characteristics and traits shared with anthropoid primates can yield insights on primate and human evolution. Most species have been discovered or promoted to full species status since the 1990s; however, lemur taxonomic classification is controversial and depends on which species concept is used.
Many lemur species remain endangered due to habitat loss and hunting. Although local traditions, such as , generally help protect lemurs and their forests, illegal logging, economic privation and political instability conspire to thwart conservation efforts. Because of these threats and their declining numbers, the International Union for Conservation of Nature (IUCN) considers lemurs to be the world's most endangered mammals, noting that up to 90% of all lemur species confront the threat of extinction in the wild within the next 20 to 25 years. Ring-tailed lemurs are an iconic flagship species. Collectively, lemurs exemplify the biodiverse fauna of Madagascar and have facilitated the emergence of eco-tourism. In addition, conservation organizations increasingly seek to implement community-based approaches to save lemur species and promote sustainability.
Etymology
Carl Linnaeus, the founder of modern binomial nomenclature, gave lemurs their name as early as 1758. With his declaration of the genus Lemur in the 10th edition of Systema Naturae, he included three species: Lemur tardigradus (the red slender loris, now known as Loris tardigradus), Lemur catta (the ring-tailed lemur), and Lemur volans (the Philippine colugo, now known as Cynocephalus volans). Although the term lemur was first intended for slender lorises, it was soon limited to the endemic Malagasy primates, which have been known as collectively "lemurs" ever since.
The name lemur is derived from the Latin term lemures, which refers to specters or ghosts that were exorcised during the Lemuria festival of ancient Rome. Linnaeus was familiar with the works of Virgil and Ovid, both of whom mentioned lemures. Seeing an analogy that fit with his naming scheme, he adapted the term "lemur" for these nocturnal primates.
It was noted in 2012 that many sources had commonly and falsely assumed that Linnaeus was referring to the ghost-like appearance, reflective eyes, and ghostly cries of Madagascar's lemurs when he selected the name. Up until then, it had also been speculated that Linnaeus may also have known that some Malagasy people believed that lemurs were the souls of their ancestors. However, both claims were discredited since according to Linnaeus' own explanation, the term lemur was selected because of the nocturnal activity and slow movements of the red slender loris:
Evolutionary history
Lemurs are primates belonging to the suborder Strepsirrhini. Like other strepsirrhine primates, such as lorises, pottos, and galagos, they share ancestral (or plesiomorphic) traits with early primates. In this regard, lemurs are popularly confused with ancestral primates; however, lemurs did not give rise to monkeys and apes (simians). Instead, they evolved independently in isolation on Madagascar. All modern strepsirrhines including lemurs are traditionally thought to have evolved from early primates known as adapiforms during the Eocene (56 to 34 mya) or Paleocene (66 to 56 mya). Adapiforms, however, lack a specialized arrangement of teeth, known as a toothcomb, which nearly all living strepsirrhines possess. A more recent hypothesis is that lemurs descended from lorisoids (loris-like) primates. This is supported by comparative studies of the cytochrome b gene and the presence of the strepsirrhine toothcomb in both groups. Instead of being the direct ancestors of lemurs, the adapiforms may have given rise to both the lemurs and lorisoids, a split that would be supported by molecular phylogenetic studies. The later split between lemurs and lorises is thought to have occurred approximately 62 to 65 mya according to molecular studies, although other genetic tests and the fossil record in Africa suggest more conservative estimates of 50 to 55 mya for this divergence. However, the oldest lemur fossils on Madagascar are actually subfossils dating to the Late Pleistocene.
Once part of the supercontinent Gondwana, the island of Madagascar has been isolated since it broke away from eastern Africa (~160 mya), Antarctica (~80–130 mya), and India (~80–90 mya). Since ancestral lemurs are thought to have originated in Africa around 62 to 65 mya, they must have crossed the Mozambique Channel, a deep channel between Africa and Madagascar with a minimum width of about 560 km (350 mi). In 1915, paleontologist William Diller Matthew noted that the mammalian biodiversity on Madagascar (including lemurs) can only be accounted for by random rafting events, where very small populations rafted from nearby Africa on tangled mats of vegetation, which get flushed out to sea from major rivers. This form of biological dispersal can occur randomly over millions of years. In the 1940s, American paleontologist George Gaylord Simpson coined the term "sweepstakes hypothesis" for such random events. Rafting has since been the most accepted explanation for the lemur colonization of Madagascar, but until recently, this trip was thought to be very unlikely because strong ocean currents flow away from the island. In , a report demonstrated that around 60 mya both Madagascar and Africa were 1,650 km (1,030 mi) south of their present-day positions, placing them in a different ocean gyre, producing currents that ran counter to what they are today. The ocean currents were shown to be even stronger than today, which would have pushed a raft along faster, shortening the trip to 30 days or less—short enough for a small mammal to survive easily. As the continental plates drifted northward, the currents gradually changed, and by 20 mya the window for oceanic dispersal had closed, effectively isolating the lemurs and the rest of the terrestrial Malagasy fauna from mainland Africa. Isolated on Madagascar with only a limited number of mammalian competitors, the lemurs did not have to compete with other evolving arboreal mammalian groups, such as squirrels. They were also spared from having to compete with monkeys, which evolved later. The intelligence, aggression, and deceptiveness of monkeys gave them an advantage over other primates in exploiting the environment.
Distribution and diversity
Lemurs have adapted to fill many open ecological niches since making their way to Madagascar. Their diversity in both behavior and morphology (outward appearance) rivals that of the monkeys and apes found elsewhere in the world. Ranging in size from the 30 g (1.1 oz) Madame Berthe's mouse lemur, the world's smallest primate, to the recently extinct 160–200 kg (350–440 lb) Archaeoindris fontoynonti, lemurs evolved diverse forms of locomotion, varying levels of social complexity, and unique adaptations to the local climate.
Lemurs lack any shared traits that make them stand out from all other primates. Different types of lemurs have evolved unique combinations of unusual traits to cope with Madagascar's harsh, seasonal climate. These traits can include seasonal fat storage, hypometabolism (including torpor and hibernation), small group sizes, low encephalization (relative brain size), cathemerality (activity both day and night), and strict breeding seasons. Extreme resource limitations and seasonal breeding are also thought to have given rise to three other relatively common lemur traits: female social dominance, sexual monomorphism, and male–male competition for mates involving low levels of agonism, such as sperm competition.
Before the arrival of humans roughly 1500 to 2000 years ago, lemurs were found all across the island. However, early settlers quickly converted the forests to rice paddies and grassland through slash-and-burn agriculture (known locally as tavy), restricting lemurs to approximately 10% of the island's area, ~60,000 km2 (23,000 sq mi). Today, the diversity and complexity of lemur communities increases with floral diversity and precipitation and is highest in the rainforests of the east coast. Despite their adaptations for weathering extreme adversity, habitat destruction and hunting have resulted in lemur populations declining sharply, and their diversity has diminished, with the recent extinction of at least 17 species in eight genera, known collectively as the subfossil lemurs. Most of the approximately 100 species and subspecies of lemur are either threatened or endangered. Unless trends change, extinctions are likely to continue.
Until recently, giant lemurs existed on Madagascar. Now represented only by recent or subfossil remains, they were modern forms that were once part of the rich lemur diversity that has evolved in isolation. Some of their adaptations were unlike those seen in their living relatives. All 17 extinct lemurs were larger than the extant (living) forms, some weighing as much as 200 kg (440 lb), and are thought to have been active during the day. Not only were they unlike the living lemurs in both size and appearance, they also filled ecological niches that either no longer exist or are now left unoccupied. Large parts of Madagascar, which are now devoid of forests and lemurs, once hosted diverse primate communities that included more than 20 lemur species covering the full range of lemur sizes.
Taxonomic classification and phylogeny
From a taxonomic standpoint, the term "lemur" originally referred to the genus Lemur, which currently contains only the ring-tailed lemur. The term is now used in the colloquial sense in reference to all Malagasy primates.
Lemur taxonomy is controversial, and not all experts agree, particularly with the recent increase in the number of recognized species. According to Russell Mittermeier, the president of Conservation International (CI), taxonomist Colin Groves, and others, there are nearly 100 recognized species or subspecies of extant (or living) lemur, divided into five families and 15 genera. Because genetic data indicates that the recently extinct subfossil lemurs were closely related to living lemurs, an additional three families, eight genera, and 17 species can be included in the total. In contrast, other experts have labeled this as taxonomic inflation, instead preferring a total closer to 50 species.
The classification of lemurs within the suborder Strepsirrhini is equally controversial, although most experts agree on the same phylogenetic tree. In one taxonomy, the infraorder Lemuriformes contains all living strepsirrhines in two superfamilies, Lemuroidea for all lemurs and Lorisoidea for the lorisoids (lorisids and galagos). Alternatively, the lorisoids are sometimes placed in their own infraorder, Lorisiformes, separate from the lemurs. In another taxonomy published by Colin Groves, the aye-aye was placed in its own infraorder, Chiromyiformes, while the rest of the lemurs were placed in Lemuriformes and the lorisoids in Lorisiformes.
Although it is generally agreed that the aye-aye is the most basal member of the lemur clade, the relationship between the other four families is less clear since they diverged during a narrow 10 to 12 million-year window between the Late Eocene (42 mya) and into the Oligocene (30 mya). The two main competing hypotheses are shown in the adjacent image.
Lemur taxonomy has changed significantly since the first taxonomic classification of lemurs by Carl Linnaeus in 1758. One of the greatest challenges has been the classification of the aye-aye, which has been a topic of debate up until very recently. Until Richard Owen published a definitive anatomical study in 1866, early naturalists were uncertain whether the aye-aye (genus Daubentonia) was a primate, rodent, or marsupial. However, the placement of the aye-aye within the order Primates remained problematic until very recently. Based on its anatomy, researchers have found support for classifying the genus Daubentonia as a specialized indriid, a sister group to all strepsirrhines, and as an indeterminate taxon within the order Primates. Molecular tests have now shown Daubentoniidae is basal to all Lemuriformes, and in 2008, Russell Mittermeier, Colin Groves, and others ignored addressing higher-level taxonomy by defining lemurs as monophyletic and containing five living families, including Daubentoniidae.
Relationships among lemur families have also proven to be problematic and have yet to be definitively resolved. To further complicate the issue, several Paleogene fossil primates from outside Madagascar, such as Bugtilemur, have been classified as lemurs. However, scientific consensus does not accept these assignments based on genetic evidence, and therefore it is generally accepted that the Malagasy primates are monophyletic. Another area of contention is the relationship between the sportive lemurs and the extinct koala lemurs (Megaladapidae). Formerly grouped in the same family due to similarities in dentition, they are no longer considered to be closely related due to genetic studies.
More taxonomic changes have occurred at the genus level, although these revisions have proven more conclusive, often supported by genetic and molecular analysis. The most noticeable revisions included the gradual split of a broadly defined genus Lemur into separate genera for the ring-tailed lemur, ruffed lemurs, and brown lemurs due to a host of morphological differences.
Due to several taxonomic revisions by Russell Mittermeier, Colin Groves, and others, the number of recognized lemur species has grown from 33 species and subspecies in 1994 to approximately 100 in 2008. With continuing cytogenetic and molecular genetic research, as well as ongoing field studies, particularly with cryptic species such as mouse lemurs, the number of recognized lemur species is likely to keep growing. However, the rapid increase in the number of recognized species has had its critics among taxonomists and lemur researchers. Since classifications ultimately depend on the species concept used, conservationists often favor definitions that result in the splitting of genetically distinct populations into separate species to gain added environmental protection. Others favor a more thorough analysis.
Anatomy and physiology
Lemurs vary greatly in size. They include the smallest primates in the world and, until recently, also included some of the largest. They currently range in size from about 30 g (1.1 oz) for Madame Berthe's mouse lemur (Microcebus berthae) up to 7–9 kg (15–20 lb) for the indri (Indri indri) and diademed sifaka (Propithecus diadema). One recently extinct species rivaled the gorilla in size, at 160–200 kg (350–440 lb) for Archaeoindris fontoynonti.
Like all primates, lemurs have five divergent digits with nails (in most cases) on their hands and feet. Most lemurs possess a laterally compressed, elongated nail, called a toilet-claw, on the second toe and use it for scratching and grooming. In addition to the toilet-claw, lemurs share a variety of other traits with other strepsirrhine primates, which include a rhinarium (or "wet nose"); a fully functional vomeronasal organ, which detects pheromones; a postorbital bar and the lack of postorbital closure (a wall of thin bone behind the eye); orbits (bony sockets that enclose the eye) that are not fully facing forward; left and right mandible (lower jaw) bones that are not fully fused; and a small brain-to-body mass ratio.
Additional traits shared with other prosimian primates (strepsirrhine primates and tarsiers) include a bicornuate (two-horned) uterus and epitheliochorial placentation. Because their thumbs are only pseudo-opposable, making their movement less independent of the other fingers, their hands are less than perfect at grasping and manipulating objects. On their feet, they have a widely abducted hallux (first toe) which facilitates the grasping of tree limbs. A common misconception is that lemurs have a prehensile tail, a trait found only in New World monkeys, particularly atelids, among primates. Lemurs also rely heavily on their sense of smell, a trait shared with most other mammals and early primates, but not with the visually oriented higher primates. This sense of smell is important in terms of marking territory as well as provide an indication of whether or not another lemur is a viable breeding partner.
Lemurs are a diverse group of primates in terms of morphology and physiology. Some lemurs, such as the sportive lemurs and indriids, have longer hind limbs than forelimbs, making them excellent leapers. Indriids also have a specialized digestive system for folivory, exhibiting enlarged salivary glands, a spacious stomach, and an elongated caecum (lower gut) that facilitates fermentation. The hairy-eared dwarf lemur (Allocebus trichotis) reportedly has a very long tongue, allowing it to feed on nectar. Likewise, the red-bellied lemur (Eulemur rubriventer) has a feathery brush-shaped tongue, also uniquely adapted to feed on nectar and pollen. The aye-aye has evolved some traits that are unique among primates, making it stand out among the lemurs. Such traits include continuously growing, rodent-like front teeth for gnawing through wood and hard seeds; a highly mobile, filiform (filament-shaped) middle finger for extracting food from tiny holes; large, bat-like ears for detecting hollow spaces within trees; and use of self-generated acoustical cues to forage.
Lemurs are unusual since they have great variability in their social structure, yet generally lack sexual dimorphism in size and canine tooth morphology. However, some species tend towards having larger females, and two species of true lemur (genus Eulemur), the gray-headed lemur (E. albocollaris) and the red lemur (E. rufus), exhibit size differences in canine teeth. True lemurs show sexual dichromatism (sexual differences in fur coloration), but the difference between the genders varies from strikingly obvious, as in the blue-eyed black lemur (E. macaco), to nearly imperceptible in the case of the common brown lemur (E. fulvus).
Crypsis, or the inability of humans to visually distinguish between two or more distinct species, has recently been discovered among lemurs, particularly within the sportive lemurs (Lepilemur) and mouse lemurs (Microcebus). With sportive lemurs, subspecies were traditionally defined based on slight morphological differences, but new genetic evidence has supported giving full species status to these regional populations. In the case of mouse lemurs, the gray mouse lemur (M. murinus), golden-brown mouse lemur (M. ravelobensis), and Goodman's mouse lemur (M. lehilahytsara) were considered the same species until recently, when genetic tests identified them as cryptic species.
Dentition
The lemur dentition is heterodont (having multiple tooth morphologies) and derives from an ancestral primate permanent dentition of . Indriids, sportive lemurs, the aye-aye, and the extinct sloth lemurs, monkey lemurs, and koala lemurs have reduced dentitions, having lost incisors, canines, or premolars. The ancestral deciduous dentition is , but young indriids, aye-ayes, koala lemurs, sloth lemurs, and probably monkey lemurs have fewer deciduous teeth.
There are also noticeable differences in dental morphology and tooth topography between lemurs. Indri, for instance, have teeth that are perfectly adapted for shearing leaves and crushing seeds. In the toothcomb of most lemurs, the bottom incisors and canine teeth are procumbent (face forward rather than up) and finely spaced, thus providing a tool for either grooming or feeding. For instance, indri use their toothcomb not only for grooming, but also to pry out the large seeds from the tough epicarp of Beilschmiedia fruits, while fork-marked lemurs use their relatively long toothcomb to cut through tree bark to induce the flow of tree sap. The toothcomb is kept clean by the sublingua or "under-tongue", a specialized structure that acts like a toothbrush to remove hair and other debris. The sublingua extends below the tip of the tongue and is tipped with keratinized, serrated points that rake between the front teeth.
Only the aye-aye, the extinct giant aye-aye, and the largest of the extinct giant sloth lemurs lack a functional strepsirrhine toothcomb. In the case of the aye-aye, the morphology of the deciduous incisors, which are lost shortly after birth, indicates that its ancestors had a toothcomb. These milk teeth are lost shortly after birth and are replaced by open-rooted, continually growing (hypselodont) incisors.
The toothcomb in lemurs normally consists of six teeth (four incisors and two canines), although indriids, monkey lemurs, and some sloth lemurs only have a four-tooth toothcomb due to the loss of either a pair of canines or incisors. Because the lower canine is either included in the toothcomb or lost, the lower dentition can be difficult to read, especially since the first premolar (P2) is often shaped like a canine (caniniform) to fill the canine's role. In folivorous (leaf-eating) lemurs, except for indriids, the upper incisors are greatly reduced or absent. Used together with the toothcomb on the mandible (lower jaw), this complex is reminiscent of an ungulate browsing pad.
Lemurs are unusual among primates for their rapid dental development, particularly among the largest species. For example, indriids have relatively slow body growth but extremely fast tooth formation and eruption. By contrast, anthropoid primates exhibit slower dental development with increased size and slower morphological development. Lemurs are also dentally precocious at birth, and have their full permanent dentition at weaning.
Lemurs generally have thin tooth enamel compared to anthropoid primates. This may result in extra wear and breakage to the anterior (front) teeth due to heavy use in grooming, feeding, and fighting. Little other dental health information is available for lemurs, except that wild ring-tailed lemurs at Berenty Private Reserve occasionally exhibit abscessed maxillary canines (seen as open wounds on the muzzle) and tooth decay, possibly due to the consumption of non-native foods.
Senses
The sense of smell, or olfaction, is highly important to lemurs and is frequently used in communication. Lemurs have long snouts (compared to the short snouts of haplorrhines) that are traditionally thought to position the nose for better sifting of smells, although long snouts do not necessarily translate into high olfactory acuity since it is not the relative size of the nasal cavity that correlates with smell, but the density of olfactory receptors. Instead, the long snouts may facilitate better chewing.
The wet nose, or rhinarium, is a trait shared with other strepsirrhines and many other mammals, but not with haplorrhine primates. Although it is claimed to enhance the sense of smell, it is actually a touch-based sense organ that connects with a well-developed vomeronasal organ (VNO). Since pheromones are usually large, non-volatile molecules, the rhinarium is used to touch a scent-marked object and transfer the pheromone molecules down the philtrum (the nasal mid-line cleft) to the VNO via the nasopalatine ducts that travel through the incisive foramen of the hard palate.
To communicate with smell, which is useful at night, lemurs will scent mark with urine as well as scent glands located on the wrists, inside elbow, genital regions, or the neck. The scrotal skin of most male lemurs has scent glands. Ruffed lemurs (genus Varecia) and male sifakas have a gland at the base of their neck, while the greater bamboo lemur (Prolemur simus) and the ring-tailed lemur have glands inside the upper arms near the axilla. Male ring-tailed lemurs also have scent glands on the inside of their forearms, adjacent to a thornlike spur, which they use to gouge, and simultaneously, scent-mark tree branches. They will also wipe their tails between their forearms and then engage in "stink fights" by waving their tail at their opponents.
Lemurs (and strepsirrhines in general) are considered to be less visually oriented than the higher primates, since they rely so heavily on their sense of smell and pheromone detection. The fovea on the retina, which yields higher visual acuity, is not well-developed. The postorbital septum (or bony closure behind the eye) in haplorrhine primates is thought to stabilize the eye slightly, allowing for the evolution of the fovea. With only a postorbital bar, lemurs have been unable to develop a fovea. Therefore, regardless of their activity pattern (nocturnal, cathemeral, or diurnal), lemurs exhibit low visual acuity and high retinal summation. Lemurs can see a wider visual field, however, than anthropoid primates due to a slight difference in the angle between the eyes, as shown in the following table:
Although they lack a fovea, some diurnal lemurs have a cone-rich, although less clustered, area centralis. This area centralis has a high rod-to-cone cell ratio in many diurnal species studied thus far, whereas diurnal anthropoids have no rod cells in their fovea. Once again, this suggests lower visual acuity in lemurs than in anthropoids. Furthermore, the rod-to-cone cell ratio can be variable even among diurnal species. For instance, Verreaux's sifaka (Propithecus verreauxi) and the indri (Indri indri) have only a few large cones scattered along their predominantly rod-dominated retina. The eyes of the ring-tailed lemur contain one cone to five rods. Nocturnal lemurs such as mouse lemurs and dwarf lemurs, on the other hand, have retinas made up entirely of rod cells.
Since cone cells make color vision possible, the high prevalence of rod cells in lemur eyes suggest they have not evolved color vision. The most studied lemur, the ring-tailed lemur, has been shown to have blue-yellow vision, but lacks the ability to distinguish red and green hues. Due to polymorphism in opsin genes, which code for color receptivity, trichromatic vision may rarely occur in females of a few lemur species, such as Coquerel's sifaka (Propithecus coquereli) and the red ruffed lemur (Varecia rubra). Most lemurs, therefore, are either monochromats or dichromats.
Most lemurs have retained the tapetum lucidum, a reflective layer of tissue in the eye, which is found in many vertebrates. This trait is absent in haplorrhine primates, and its presence further limits the visual acuity in lemurs. The strepsirrhine choroidal tapetum is unique among mammals because it is made up of crystalline riboflavin, and the resulting optical scattering is what limits visual acuity. Although the tapetum is considered to be ubiquitous in lemurs, there appear to be exceptions among true lemurs, such as the black lemur and the common brown lemur, as well as the ruffed lemurs. Since the riboflavins in the tapetum have a tendency to dissolve and vanish when processed for histological investigation, however, the exceptions are still debatable.
Lemurs also have a third eyelid known as a nictitating membrane, whereas most other primates have a lesser developed plica semilunaris. The nictitating membrane keeps the cornea moist and clean by sweeping across the eye.
Metabolism
Lemurs have low basal metabolic rates (BMR), which helps them to conserve energy during the dry season, when water and food are scarce. They can optimize their energy use by lowering their metabolic rate to 20% below the values predicted for mammals of similar body mass. The red-tailed sportive lemur (Lepilemur ruficaudatus), for instance, reportedly has one of the lowest metabolic rates among mammals. Its low metabolic rate may be linked to its generally folivorous diet and relatively small body mass. Lemurs exhibit behavioral adaptations to complement this trait, including sunning behaviors, hunched sitting, group huddling, and nest sharing, in order to reduce heat loss and conserve energy. Dwarf lemurs and mouse lemurs exhibit seasonal cycles of dormancy to conserve energy. Before dry season, they will accumulate fat in white adipose tissue located at the base of the tail and hind legs, doubling their weight. At the end of the dry season, their body mass may fall to half of what it was prior to the dry season. Lemurs that do not experience states of dormancy are also able to shut down aspects of their metabolism for energy conservation.
Behaviour
Lemur behaviour is as variable as lemur morphology. Differences in diet, social systems, activity patterns, locomotion, communication, predator avoidance tactics, breeding systems, and intelligence levels help define lemur taxa and set individual species apart from the rest. Although trends frequently distinguish the smaller, nocturnal lemurs from the larger, diurnal lemurs, there are often exceptions that help exemplify the unique and diverse nature of these Malagasy primates.
Diet
Lemur diets are highly variable and demonstrate a high degree of plasticity, although general trends suggest that the smallest species primarily consume fruit and insects (omnivory), while the larger species are more herbivorous, consuming mostly plant material. As with all primates, hungry lemurs might eat anything that is edible, whether or not the item is one of their preferred foods. For instance, the ring-tailed lemur eats insects and small vertebrates when necessary and as a result it is commonly viewed as an opportunistic omnivore. Coquerel's giant mouse lemur (Mirza coquereli) is mostly frugivorous, but will consume insect secretions during the dry season.
A common assumption in mammalogy is that small mammals cannot subsist entirely on plant material and must have a high-calorie diet in order to survive. As a result, it was thought that the diet of tiny primates must be high in protein-containing insects (insectivory). Research has shown, however, that mouse lemurs, the smallest living primates, consume more fruit than insects, contradicting the popular hypothesis.
Plant material makes up the majority of most lemur diets. Members of at least 109 of all known plant families in Madagascar (55%) are exploited by lemurs. Since lemurs are primarily arboreal, most of these exploited species are woody plants, including trees, shrubs, or lianas. Only the ring-tailed lemur, the bamboo lemurs (genus Hapalemur), and the black-and-white ruffed lemur (Varecia variegata) are known to consume herbs. While Madagascar is rich in fern diversity, these plants are rarely eaten by lemurs. One possible reason for this is that ferns lack flowers, fruits, and seeds—common food items in lemur diets. They also occur close to the ground, while lemurs spend most of their time in the trees. Lastly, ferns have an unpleasant taste due to the high content of tannins in their fronds. Likewise, mangroves appear to be rarely exploited by lemurs due to their high tannin content. Some lemurs appear to have evolved responses against common plant defenses, however, such as tannins and alkaloids. The golden bamboo lemur (Hapalemur aureus), for instance, eats giant bamboo (Cathariostachys madagascariensis), which contains high levels of cyanide. This lemur can consume twelve times the typically lethal dose for most mammals on a daily basis; the physiological mechanisms that protect it from cyanide poisoning are unknown. At the Duke Lemur Center (DLC) in the United States, lemurs that roam the outdoor enclosures have been observed eating poison ivy (Taxicodendron radicans), yet have shown no ill effects.
Many of the larger lemur species consume leaves (folivory), particularly the indriids. However, some smaller lemurs such as sportive lemurs (genus Lepilemur) and woolly lemurs (genus Avahi) also primarily eat leaves, making them the smallest primates that do so. The smallest of the lemurs generally do not eat much leaf matter. Collectively, lemurs have been documented consuming leaves from at least 82 native plant families and 15 alien plant families. Lemurs tend to be selective in their consumption of the part of the leaf or shoot as well as its age. Often, young leaves are preferred over mature leaves.
Many lemurs that eat leaves tend to do so during times of fruit scarcity, sometimes suffering weight loss as a result. Most lemur species, including most of the smallest lemurs and excluding some of the indriids, predominantly eat fruit (frugivory) when available. Collectively, lemurs have been documented consuming fruit from at least 86 native plant families and 15 alien plant families. As with most tropical fruit eaters, the lemur diet is dominated by fruit from Ficus (fig) species. In many anthropoid primates, fruit is a primary source of vitamin C, but unlike anthropoid primates, lemurs (and all strepsirrhines) can synthesize their own vitamin C. Historically, captive lemur diets high in vitamin C-rich fruits have been thought to cause hemosiderosis, a type of iron overload disorder, since vitamin C increases iron absorption. Although lemurs in captivity have been shown to be prone to hemosiderosis, the frequency of the disease varies across institutions and may depend on the diet, husbandry protocols, and genetic stock. Assumptions about the problem need to be tested separately for each species. The ring-tailed lemur, for instance, seems to be less prone to the disorder than other lemur species.
Only eight species of lemur are known to be seed predators (granivores), but this may be under-reported since most observations only report fruit consumption and do not investigate whether the seeds are consumed as well. These lemurs include some indriids, such as the diademed sifaka (Propithecus diadema), the golden-crowned sifaka (Propithecus tattersalli), the indri, and the aye-aye. The aye-aye, which specializes in structurally defended resources, can chew through Canarium seeds, which are harder than the seeds that New World monkeys are known to break open. At least 36 genera from 23 families of plants are targeted by lemur seed predators.
Inflorescences (clusters of flowers) of at least 60 plant families are eaten by lemurs ranging in size from the tiny mouse lemurs to the relatively large ruffed lemurs. If the flowers are not exploited, sometimes the nectar is consumed (nectarivory) along with the pollen (palynivory). At least 24 native species from 17 plant families are targeted for nectar or pollen consumption.
Bark and plant exudates such as tree sap are consumed by a few lemur species. The exploitation of exudates has been reported in 18 plant species and only in the dry regions in the south and west of Madagascar. Only the Masoala fork-marked lemur (Phaner furcifer) and Coquerel's giant mouse lemur regularly consume tree sap. Bark has never been reported as an important food item in lemur diets, but at least four species eat it: the aye-aye, the red-tailed sportive lemur (Lepilemur ruficaudatus), the common brown lemur (Eulemur fulvus), and Verreaux's sifaka (Propithecus verreauxi). Most bark feeding is directly linked to exudate feeding, except for the aye-aye's bark feeding on Afzelia bijuga (genus Afzelia) at Nosy Mangabe in the northeast.
Soil consumption (geophagy) has also been reported and likely helps with digestion, provides minerals and salts, and helps absorb toxins. Sifakas have been observed eating soil from termite mounds, possibly adding beneficial intestinal flora to aid the digestion of cellulose from their folivorous diet.
Social systems
Lemurs are social and live in groups that usually include fewer than 15 individuals. Observed social organization patterns include "solitary but social", "fission-fusion", "pair bonds", and "multi-male group". Nocturnal lemurs are mostly solitary but social, foraging alone at night but often nesting in groups during the day. The degree of socialization varies by species, gender, location, and season. In many nocturnal species, for instance, the females, along with their young, will share nests with other females and possibly one male, whose larger home range happens to overlap one or more female nesting groups. In sportive lemurs and fork-marked lemurs, one or two females may share a home range, possibly with a male. In addition to sharing nests, they will also interact vocally or physically with their range-mate while they forage at night. Diurnal lemurs exhibit many of the social systems seen in monkeys and apes, living in relatively permanent and cohesive social groups. Multi-male groups are the most common, just as they are in most anthropoid primates. True lemurs utilize this social system, often living in groups of ten or less. Ruffed lemurs have been shown to live in fission-fusion societies, and Indri forms pair bonds.
Some lemurs exhibit female philopatry, where females stay within their natal range and the males migrate upon reaching maturity, and in other species both sexes will migrate. In some cases, female philopatry may help explain the evolution of female-bonded multi-male groups, such as those of the ring-tailed lemur, Milne-Edwards' sifaka (Propithecus edwardsi), and the Verreaux's sifaka. Their ancestors may have been more solitary, with females that lived in mother-daughter pairs (or dyads). Over time, these dyads may have allied themselves with other neighboring mother-daughter dyads in order to defend more distributed resources in a wide home range. If this is true, then multi-male groups in lemurs may differ fundamentally in their internal structure from those in catarrhine primates (Old World monkeys and apes).
The presence of female social dominance sets lemurs apart from most other primates and mammals; in most primate societies, males are dominant unless females band together to form coalitions that displace them. However, many Eulemur species are exceptions and the greater bamboo lemur (Prolemur simus) does not exhibit female dominance. When females are dominant within a group, the way they maintain dominance varies. Ring-tailed lemur males act submissively with or without signs of female aggression. Male crowned lemurs (Eulemur coronatus), on the other hand, will only act submissively when females act aggressively towards them. Female aggression is often associated with, but not limited to, feeding.
There have been many hypotheses that have attempted to explain why lemurs exhibit female social dominance while other primates with similar social structures do not, but no consensus has been reached after decades of research. The dominant view in the literature states that female dominance is an advantageous trait given the high costs of reproduction and the scarcity of resources available. Indeed, female dominance has been shown to be linked to increased maternal investment. However, when reproductive costs and extreme seasonality of resources were compared across primates, other primates demonstrated male dominance under conditions that were similar to or more challenging than those faced by lemurs. In 2008, a new hypothesis revised this model using simple game theory. It was argued that when two individuals were equally matched in fighting capacity, the one with the most need would win the conflict since it would have the most to lose. Consequently, the female, with higher resource needs for pregnancy, lactation, and maternal care, was more likely to win in resource conflicts with equally sized males. This, however, assumed monomorphism between sexes. The following year, a new hypothesis was proposed to explain monomorphism, stating that because most female lemurs are only sexually receptive for a day or two each year, males can utilize a more passive form of mate guarding: copulatory plugs, which block the female reproductive tract, preventing other males from successfully mating with her, and thus reducing the need for aggression and the evolutionary drive for sexual dimorphism.
In general, levels of agonism (or aggression) tend to correlate with relative canine height. The ring-tailed lemur has long, sharp upper canine teeth in both sexes, and it also exhibits high levels of agonism. The Indri, on the other hand, has smaller canines and exhibits lower levels of aggression. When neighboring groups of the same species defend their territories, the conflict can take the form of ritualized defense. In sifakas, these ritualized combats involve staring, growling, scent-marking, and leaping to occupy certain sections of the tree. The indri defends its home range with ritualized "singing" battles.
Like other primates, lemurs groom socially (allogroom) to ease tensions and solidify relationships. They groom in greeting, when waking up, when settling in for sleep, between mother and infant, in juvenile relations, and for sexual advances. Unlike anthropoid primates, who part the fur with the hands and pick out particles with the fingers or mouth, lemurs groom with their tongue and scraping with their toothcomb. Despite the differences in technique, lemurs groom with the same frequency and for the same reasons as anthropoids.
Activity patterns
The biological rhythm can vary from nocturnal in smaller lemurs to diurnal in most larger lemurs. Diurnality is not seen in any other living strepsirrhine. Cathemerality, where an animal is active sporadically both day and night, occurs among some of the larger lemurs. Few if any other primates exhibit this sort of activity cycle, either regularly or irregularly under changing environmental conditions. The most heavily studied cathemeral lemurs are the true lemurs. Although the mongoose lemur (E. mongoz) is the best-documented example, every species in the genus studied has shown some degree of cathemeral behavior, although night activity is often restricted by light availability and moon periodicity. This type of behavior was first documented in the 1960s in true lemur species as well as other Lemuridae species, such as ruffed lemurs and bamboo lemurs. Initially described as "crepuscular" (active at dawn and dusk), anthropologist Ian Tattersall stimulated additional research and coined the new term "cathemeral", although many non-anthropologists prefer the terms "circadian" or "diel".
In order to conserve energy and water in their highly seasonal environment, mouse lemurs and dwarf lemurs exhibit seasonal behavioral cycles of dormancy where the metabolic rate and body temperature are lowered. They are the only primates known to do so. They accumulate fat reserves in their hind legs and the base of their tail before the dry winter season, when food and water are scarce, and can exhibit daily and prolonged torpor during the dry season. Daily torpor constitutes less than 24 hours of dormancy, whereas prolonged torpor averages two weeks in duration and signals hibernation. Mouse lemurs have been observed experiencing torpor that lasts for several consecutive days, but dwarf lemurs are known to hibernate for six to eight months every year, particularly on the west coast of Madagascar.
Dwarf lemurs are the only primates known to hibernate for extended periods. Unlike other hibernating mammals from temperate regions, which have to awaken regularly for a few days, dwarf lemurs experience five months of continuous deep hibernation (May through September). Before and after this deep hibernation, there are two months (April and October) of transition, where they will forage on a limited basis to reduce demands on their fat reserves. Unlike any other hibernating mammal, the body temperature of hibernating dwarf lemurs will fluctuate with the ambient temperature rather than remaining low and stable.
Other lemurs that do not exhibit dormancy conserve energy by selecting thermoregulated microhabitats (such as tree holes), sharing nests, and reducing exposed body surfaces, such as by hunched sitting and group huddling. Also, the ring-tailed lemur, ruffed lemurs, and sifakas are commonly seen sunning, thus using solar radiation to warm their bodies instead of metabolic heat.
Locomotion
Locomotor behavior in lemurs, both living and extinct, is highly varied and its diversity exceeds that of anthropoids. Locomotor postures and behaviors have included vertical clinging and leaping (including saltatory behavior), seen in indriids and bamboo lemurs; slow (loris-like) arboreal quadrupedal locomotion, once exhibited by Mesopropithecus; fast arboreal quadrupedal locomotion, seen in true lemurs and ruffed lemurs; partially terrestrial quadrupedal locomotion, seen in the ring-tailed lemur; highly terrestrial quadrupedal locomotion, once exhibited by monkey lemurs such as Hadropithecus; and sloth-like suspensory locomotion, once exhibited by many of the sloth lemurs, such as Palaeopropithecus. The Lac Alaotra gentle lemur (Hapalemur alaotrensis) has even been reported to be a good swimmer. Sometimes these locomotor types are lumped together into two main groups of lemurs, the vertical clingers and leapers and the arboreal (and occasionally terrestrial) quadrupeds.
The jumping prowess of the indriids has been well documented and is popular among ecotourists visiting Madagascar. Using their long, powerful back legs, they catapult themselves into the air and land in an upright posture on a nearby tree, with both hands and feet tightly gripping the trunk. Indriids can leap up to 10 m (33 ft) rapidly from tree trunk to tree trunk, an ability referred to as "ricochetal leaping". Verreaux's sifaka (Propithecus verreauxi) manages to do this in the spiny forests of southern Madagascar. It is unknown how it avoids impaling its palms on the thorn-covered trunks of large plants such as Alluaudia. When distances between trees are too great, sifakas will descend to the ground and cross distances more than 100 m (330 ft) by standing upright and hopping sideways with the arms held to the side and waving up and down from chest to head height, presumably for balance. This is sometimes described as a "dance-hop".
Communication
Lemur communication can be transmitted through sound, sight, and smell (olfaction). The ring-tailed lemur, for instance, uses complex though highly stereotyped behaviors such as scent-marking and vocalizations. Visual signals are probably the least used by lemurs, since they lack many of the muscles used in common primate facial expressions. Given their poor vision, whole-body postures are probably more noticeable. However, the ring-tailed lemur has demonstrated distinct facial expressions including a threat stare, pulled back lips for submission, and pulled back ears along with flared nostrils during scent-marking. This species has also been observed using yawns as threats. Their ringed tails also communicate distance, warn off neighboring troops, and help locate troop members. Sifakas are known to exhibit an open-mouth play face as well as a submissive teeth-baring grimace used in agonistic interactions.
Olfaction is particularly important to lemurs, except for the indri, which lacks most common lemur scent glands and has a greatly reduced olfactory region in the brain. Olfaction can communicate information about age, sex, reproductive status, as well as demarcate the boundaries of a territory. It is most useful for communication between animals that rarely encounter each other. Small, nocturnal lemurs mark their territories with urine, while the larger, diurnal species use scent glands located on various parts of their anatomy. The ring-tailed lemur engages in "stink fights" by rubbing its tail across scent glands on its wrists and then flicking it at other male opponents. Some lemurs defecate in specific areas, otherwise known as latrine behavior. Although many animals exhibit this behavior, it is a rare trait among primates. Latrine behavior can represent territorial marking and aid in interspecies signaling.
Compared to other mammals, primates in general are very vocal, and lemurs are no exception. Some lemur species have extensive vocal repertoires, including the ring-tailed lemur and ruffed lemurs. Some of the most common calls among lemurs are predator alarm calls. Lemurs not only respond to alarm calls of their own species, but also alarm calls of other species and those of non-predatory birds. The ring-tailed lemur and a few other species have different calls and reactions to specific types of predators. With mating calls, it has been shown that mouse lemurs that cannot be discerned visually respond more strongly to the calls of their own species, particularly when exposed to the calls of other mouse lemurs that they would encounter normally within their home range. Lemur calls can also be very loud and carry long distances. Ruffed lemurs use several loud calls that can be heard up to 1 km (0.62 mi) away on a clear, calm day. The loudest lemur is the indri, whose calls can be heard up to 2 km (1.2 mi) or more and thus communicate more effectively the territorial boundaries over its 34 to 40 hectares (0.13 to 0.15 sq mi) home range. Both ruffed lemurs and the indri exhibit contagious calling, where one individual or group starts a loud call and others within the area join in. The song of the indri can last 45 seconds to more than 3 minutes and tends to coordinate to form a stable duet comparable to that of gibbons.
Tactile communication (touch) is mostly used by lemurs in the form of grooming, although the ring-tailed lemur also clumps together to sleep (in an order determined by rank), reaches out and touches adjacent members, and cuffs other members. Reaching out and touching another individual in this species has been shown to be a submissive behavior, done by younger or submissive animals towards older and more dominant members of the troop. Allogrooming, however, appears to occur more frequently between higher ranking individuals, a shared trait with other primate species. Unlike anthropoid primates, lemur grooming seems to be more intimate and mutual, often directly reciprocated. Anthropoids, on the other hand, use allogrooming to manage agonistic interactions. The ring-tailed lemur is known to be very tactile, spending between 5 and 11% of its time grooming.
Predator avoidance
All lemurs experience some predation pressure. Common defenses against predation include the use of alarm calls and predator mobbing, mostly among diurnal lemurs. The leaping abilities of lemurs may have evolved for predator avoidance rather than for travel, according to a study in kinematics. Nocturnal lemurs are difficult to see and track at night and decrease their visibility by foraging alone. They also try to avoid predators by using concealing sleeping locations, such as nests, tree holes, or dense vegetation, Some may also avoid areas frequented by predators by detecting the smell of their feces and alternating between multiple sleeping locations. Even torpor and hibernation states among cheirogaleids may be partly due to high levels of predation. Infants are protected while foraging by either leaving them in the nest or by stashing them in a hidden location, where the infant remains immobile in the absence of the parent.
Diurnal lemurs are visible during the day, so many live in groups, where the increased number of eyes and ears helps aid in predator detection. Diurnal lemurs use and respond to alarm calls, even those of other lemur species and non-predatory birds. The ring-tailed lemur has different calls and reactions to different classes of predators, such as predatory birds, mammals, or snakes. Some lemurs, such as the indri, use crypsis to camouflage themselves. They are often heard but difficult to see in the trees due to the dappled light, earning them the reputation of being "ghosts of the forest".
Reproduction
Except for the aye-aye and the Lac Alaotra gentle lemur, lemurs are seasonal breeders with very short mating and birth seasons influenced by the highly seasonal availability of resources in their environment. Mating season usually last less than three weeks each year, and the female vagina opens up only during a few hours or days of her most receptive time of estrus. These narrow windows for reproduction and resource availability appear to relate to their short gestation periods, rapid maturation, and low basal metabolic rates, as well as the high energy costs of reproduction for females. This may also relate to the relatively high mortality rate among adult females and the higher proportion of adult males in some lemur populations—both unusual traits among primates. In both the aye-aye and Lac Alaotra gentle lemur, birth (parturition) occurs over a six-month period.
Lemurs time their mating and birth seasons so that all weaning periods are synchronized to match the time of highest food availability. Weaning occurs either before or shortly after the eruption of the first permanent molars in lemurs. Mouse lemurs are able to fit their entire breeding cycle into the wet season, whereas larger lemurs, such as sifakas, must lactate for two months during the dry season. Infant survival in some species, such as Milne-Edwards' sifaka, has been shown to be directly impacted by both environmental conditions and the rank, age, and health of the mother. The breeding season is also affected by geographical location. For example, mouse lemurs give birth between September and October in their native habitat in the Southern Hemisphere, but from May through June in the captive settings in the Northern Hemisphere.
Scent factors heavily into lemur reproduction. Scent-marking activity escalates during the mating season. Pheromones may coordinate reproductive timing for females coming into estrus. Mating can be either monogamous or promiscuous for both males and females, and mating can include individuals from outside the group. Monogamous lemurs include the red-bellied lemur (Eulemur rubriventer) and the mongoose lemur (E. mongoz), although the mongoose lemur has been observed mating outside of its pair bond. Monogamy is most common among nocturnal species, although some exhibit scramble competition, sexual suppression of subordinates, or competitions between males that avoid direct fighting. In mouse lemurs, males utilize sperm plugs, developed enlarged testes during the mating season, and develop size dimorphism (likely due to the enlarged testes). These indicate a mating system known as scramble competition polygyny, where males cannot defend females or the resources that might attract them.
The gestation period varies within lemurs, ranging from 9 weeks in mouse lemurs and 9–10 weeks in dwarf lemurs to 18–24 weeks in other lemurs. The smaller, nocturnal lemurs, such as mouse lemurs, giant mouse lemurs, and dwarf lemurs, usually give birth to more than one infant, whereas the larger, nocturnal lemurs, such as fork-marked lemurs, sportive lemurs, and the aye-aye usually have one offspring. Dwarf and mouse lemurs have up to four offspring, but both average only two. Ruffed lemurs are the only large, diurnal lemurs to consistently give birth to two or three offspring. All other lemurs have single births. Multiple births in lemurs are normally fraternal, and are known to occur in every five to six births in species such as the ring-tailed lemur and some Eulemur.
After the offspring are born, lemurs either carry them around or stash them while foraging. When transported, the infants either cling to the mother's fur or are carried in the mouth by the scruff. In some species, such as bamboo lemurs, infants are carried by mouth until they are able to cling to their mother's fur. Species that park their offspring include nocturnal species (e.g. mouse lemurs, sportive lemurs, and dwarf lemurs), bamboo lemurs, and ruffed lemurs. In the case of the ruffed lemurs, the young are altricial and the mothers build nests for them, much like the smaller, nocturnal lemur species. Woolly lemurs are unusual for nocturnal lemurs because they live in cohesive family groups and carry their single offspring with them rather than parking them. Alloparenting (multiple or group parenting) has been reported in all lemur families except the sportive lemurs and aye-aye. Allonursing is also known to occur in several lemur groups. Even males have been observed caring for infants in species such as the red-bellied lemur, mongoose lemur, eastern lesser bamboo lemur, silky sifaka, fat-tailed dwarf lemur, and ruffed lemurs.
Yet another trait that sets most lemurs apart from anthropoid primates is their long lifespan together with their high infant mortality. Many lemurs, including the ring-tailed lemur, have adapted to a highly seasonal environment, which has affected their birthrate, maturation, and twinning rate (r-selection). This helps them to recover rapidly from a population crash. In captivity, lemurs can live twice as long as they do in the wild, benefiting from consistent nutrition that meets their dietary requirements, medical advancements, and improved understanding of their housing requirements. In 1960, it was thought that lemurs could live between 23 and 25 years. It is now known that the larger species can live for more than 30 years without showing signs of aging (senescence) and still be capable of reproduction.
Cognitive abilities and tool use
Lemurs have traditionally been regarded as being less intelligent than anthropoid primates, with monkeys and apes often described as having more cunning, guile, and deceptiveness. Many lemur species, such as sifakas and the ring-tailed lemur, have scored lower on tests designed for monkeys while performing as well as monkeys on other tests. These comparisons may not be fair since lemurs prefer to manipulate objects with their mouths (rather than their hands) and only take interest in objects when in captivity. Recent studies have shown that lemurs exhibit levels of technical intelligence on par with many other primates, although they manipulate objects less often. Tool use has not been witnessed by lemurs in the wild, although in captivity the common brown lemur and the ring-tailed lemur have been demonstrated to be able to understand and use tools.
A few lemurs have been noted to have relatively large brains. The extinct Hadropithecus was as large as a large male baboon and had a comparably sized brain, giving it the largest brain size relative to body size among all prosimians. The aye-aye also has a large brain-to-body ratio, which may indicate a higher level of intelligence. However, despite having a built-in tool in the form of its thin, elongated middle finger, which it uses to fish for insect grubs, the aye-aye has tested poorly in the use of extraneous tools.
Ecology
See above: Diet, Metabolism, Activity patterns, and Locomotion
Madagascar not only contains two radically different climatic zones, the rainforests of the east and the dry regions of the west, but also swings from extended drought to cyclone-generated floods. These climatic and geographical challenges, along with poor soils, low plant productivity, wide ranges of ecosystem complexity, and a lack of regularly fruiting trees (such as fig trees) have driven the evolution of lemurs' immense morphological and behavioral diversity. Their survival has required the ability to endure the persistent extremes, not yearly averages.
Lemurs have either presently or formerly filled the ecological niches normally occupied by monkeys, squirrels, woodpeckers, and grazing ungulates. With the diversity of adaptations for specific ecological niches, habitat selection among lemur families and some genera is often very specific, thus minimizing competition. In nocturnal lemurs from the more seasonal forests in the west, up to five species can coexist during the wet season due to high food abundance. However, to endure the extreme dry season, three of the five species utilize different dietary patterns and their underlying physiological traits to allow them to coexist: fork-marked lemurs feed on tree gum, sportive lemurs feed on leaves, and giant mouse lemurs sometimes feed on insect secretions. The other two species, the gray mouse lemur and the fat-tailed dwarf lemur (Cheirogaleus medius), avoid competition through reduced activity. The gray mouse lemur uses bouts of torpor, while the fat-tailed dwarf lemur hibernates completely. Similarly, on the east coast entire genera focus on specific food to avoid too much niche overlap. True lemurs and ruffed lemurs are frugivorous, indriids are folivorous, and bamboo lemurs specialize in bamboo and other grasses. Once again, seasonal dietary differences as well as subtle differences in substrate preferences, forest strata used, activity cycle, and social organization enable lemur species to coexist, although this time the species are more closely related and have similar niches. A classic example involves resource partitioning between three species of bamboo lemur that live in close proximity in small forested areas: the golden bamboo lemur, the greater bamboo lemur, and the eastern lesser bamboo lemur (Hapalemur griseus). Each utilizes either different species of bamboo, different parts of the plant, or different layers in the forest. Nutrient and toxin content (such as cyanide) help regulate food selection, though seasonal food preferences are also known to play a role.
Dietary regimes of lemurs include folivory, frugivory, and omnivory, with some being highly adaptable while others specialize on foods such as plant exudates (tree gum) and bamboo. In some cases, lemur feeding patterns directly benefit the native plant life. When lemurs exploit nectar, they may act as pollinators as long as the functional parts of the flower are not damaged. In fact, several unrelated Malagasy flowering plants demonstrate lemur-specific pollination traits, and studies indicate that some diurnal species, such as the red-bellied lemur and the ruffed lemurs, act as major pollinators. Two examples of plant species that rely on lemurs for pollination include traveller's palm (Ravenala madagascariensis) and a species of legume-like liana, Strongylodon cravieniae. Seed dispersal is another service lemurs provide. After passing through the lemur gut, tree and vine seeds exhibit lower mortality and germinate faster. Latrine behavior exhibited by some lemurs may help improve soil quality and facilitate seed dispersal. Because of their importance in maintaining a healthy forest, frugivorous lemurs may qualify as keystone mutualists.
All lemurs, particularly the smaller species, are affected by predation and they are important prey items for predators. Humans are the most significant predator of diurnal lemurs, despite taboos that occasionally forbid the hunting and eating of certain lemur species. Other predators include native euplerids, such as the fossa, feral cats, domestic dogs, snakes, diurnal birds of prey, and crocodiles. Extinct giant eagles, including one or two species from the genus Aquila and the giant Malagasy crowned eagle (Stephanoaetus mahery), as well as the giant fossa (Cryptoprocta spelea), previously also preyed on lemurs, perhaps including the giant subfossil lemurs or their subadult offspring. The existence of these extinct giants suggests that predator-prey interactions involving lemurs were more complex than they are today. Today, predator size restricts owls to the smaller lemurs, usually 100 g (3.5 oz) or less, while the larger lemurs fall victim to the larger diurnal birds of prey, such as the Madagascar harrier-hawk (Polyboroides radiatus) and the Madagascar buzzard (Buteo brachypterus).
Research
Similarities that lemurs share with anthropoid primates, such as diet and social organization, along with their own unique traits, have made lemurs the most heavily studied of all mammal groups on Madagascar. Research often focuses on the link between ecology and social organization, but also on their behavior and morphophysiology (the study of anatomy in relation to function). Studies of their life-history traits, behavior and ecology help understanding of primate evolution, since they are thought to share similarities with ancestral primates.
Lemurs have been the focus of monographic series, action plans, field guides, and classic works in ethology. However, few species have been thoroughly studied to date, and most research has been preliminary and restricted to a single locality. Only recently have numerous scientific papers been published to explain the basic aspects of behavior and ecology of poorly known species. Field studies have given insights on population dynamics and evolutionary ecology of most genera and many species. Long-term research focused on identified individuals is in its infancy and has only been started for a few populations. However, learning opportunities are dwindling as habitat destruction and other factors threaten the existence of lemur populations across the island.
Lemurs are mentioned in sailors' voyage logs as far back as 1608 and in 1658 that at least seven lemur species were described in detail by the French merchant, Étienne de Flacourt, who may also have been the only westerner to see and chronicle the existence of a giant (now extinct) lemur, which he called the tretretretre. Around 1703 merchants and sailors began bringing lemurs back to Europe, at which time James Petiver, an apothecary in London, described and illustrated the mongoose lemur. Starting in 1751, the London illustrator George Edwards began describing and illustrating some lemur species, of which a few were included in various editions of Systema Naturae by Carl Linnaeus. In the 1760s and 1770s, French naturalists Georges-Louis Leclerc, Comte de Buffon and Louis-Jean-Marie Daubenton began describing the anatomy of several lemur species. The first traveling naturalist to comment on lemurs was Philibert Commerçon in 1771, although it was Pierre Sonnerat who recorded a greater variety of lemur species during his travels.
During the 19th century, there was an explosion of new lemur descriptions and names, which later took decades to sort out. During this time, professional collectors gathered specimens for museums, menageries, and cabinets. Some of the major collectors were Johann Maria Hildebrandt and Charles Immanuel Forsyth Major. From these collections, as well as increasing observations of lemurs in their natural habitats, museum systematists including Albert Günther and John Edward Gray continued to contribute new names for new lemur species. However, the most notable contributions from this century includes the work of Alfred Grandidier, a naturalist and explorer who devoted himself to the study of Madagascar's natural history and local people. With the help of Alphonse Milne-Edwards, most of the diurnal lemurs were illustrated at this time. However, lemur taxonomic nomenclature took its modern form in the 1920s and 1930s, being standardized by Ernst Schwarz in 1931.
Although lemur taxonomy had developed, it was not until the 1950s and 1960s that the in-situ (or on-site) study of lemur behavior and ecology began to blossom. Jean-Jacques Petter and Arlette Petter-Rousseaux toured Madagascar in 1956 and 1957, surveying many of its lemur species and making important observations about their social groupings and reproduction. In 1960, the year of Madagascar's independence, David Attenborough introduced lemurs to the West with a commercial film. Under the guidance of John Buettner-Janusch, who founded the Duke Lemur Center in 1966, Alison Jolly traveled to Madagascar in 1962 to study the diet and social behavior of the ring-tailed lemur and Verreaux's sifaka at Berenty Private Reserve. The Petters and Jolly spawned a new era of interest in lemur ecology and behavior and were shortly followed by anthropologists such as Alison Richard, Robert Sussman, Ian Tattersall, and many others. Following the political turmoil of the mid-1970s and Madagascar's revolution, field studies resumed in the 1980s, thanks in part to the renewed involvement of the Duke Lemur Center under the direction of Elwyn L. Simons and the conservation efforts of Patricia Wright. In the decades that followed, huge strides have been made in lemur studies and many new species have been discovered.
Ex situ research (or off-site research) is also popular among researchers looking to answer questions that are difficult to test in the field. For example, efforts to sequence the genome of the gray mouse lemur will help researchers understand which genetic traits set primates apart from other mammals and will ultimately help understand what genomic traits set humans apart from other primates. One of the foremost lemur research facilities is the Duke Lemur Center (DLC) in Durham, North Carolina. It maintains the largest captive lemur population outside of Madagascar, which it maintains for non-invasive research and captive breeding. Many important research projects have been carried out there, including studies of lemur vocalizations, basic locomotor research, the kinematics of bipedalism, the effects of social complexity transitive reasoning, and cognition studies involving a lemur's ability to organize and retrieve sequences from memory. Other facilities, such as the Lemur Conservation Foundation, located near Myakka City, Florida, have also hosted research projects, such as one that looked at lemurs' ability to preferentially select tools based on functional qualities.
Conservation
Lemurs are threatened by a host of environmental problems, including deforestation, hunting for bushmeat, live capture for the exotic pet trade, and climate change. All species are listed by CITES on Appendix I, which prohibits trade of specimens or parts, except for scientific purposes. As of 2005, the International Union for Conservation of Nature (IUCN) listed 16% of all lemur species as critically endangered, 23% as endangered, 25% as vulnerable, 28% as "data deficient", and only 8% as least concern. Over the next five years, at least 28 species were newly identified, none of which have had their conservation status assessed. Many are likely to be considered threatened since the new lemur species that have been described recently are typically confined to small regions. Given the rate of continued habitat destruction, undiscovered species could go extinct before being identified. Since the arrival of humans on the island approximately 2000 years ago, all endemic Malagasy vertebrates over 10 kg (22 lb) have disappeared, including 17 species, 8 genera, and 3 families of lemurs. The IUCN Species Survival Commission (IUCN/SSC), the International Primatological Society (IPS), and Conservation International (CI) have included as many as five lemurs in their biennial "Top 25 Most Endangered Primates". The 2008–2010 list includes the greater bamboo lemur, gray-headed lemur (Eulemur cinereiceps), blue-eyed black lemur (Eulemur flavifrons), northern sportive lemur (Lepilemur septentrionalis), and silky sifaka. In 2012, an assessment by the Primate Specialist Group of the International Union for Conservation of Nature (IUCN) concluded that 90% of the then 103 described species of lemur should be listed as threatened on the IUCN Red List, making lemurs the most endangered group of mammals. The IUCN reiterated its concern in 2013, noting that 90% of all lemur species could be extinct within 20 to 25 years unless a US$7 million 3-year conservation plan aimed at helping local communities can be implemented.
Madagascar is one of the poorest countries in the world, with a high population growth rate of 2.5% per year and nearly 70% of the population living in poverty. The country is also burdened with high levels of debt and limited resources. These socioeconomic issues have complicated conservation efforts, even though the island of Madagascar has been recognized by IUCN/SSC as a critical primate region for over 20 years. Due to its relatively small land area—587,045 km2 (226,659 sq mi)—compared to other high-priority biodiversity regions and its high levels of endemism, the country is considered one of the world's most important biodiversity hotspots, with lemur conservation being a high priority. Despite the added emphasis for conservation, there is no indication that the extinctions that began with the arrival of humans have come to an end.
Threats in the wild
The greatest concern facing lemur populations is habitat destruction and degradation. Deforestation takes the form of local subsistence use, such as slash and burn agriculture (referred to as tavy in Malagasy), the creation of pasture for cattle through burning, and legal and illegal gathering of wood for firewood or charcoal production; commercial mining; and the illegal logging of precious hardwoods for foreign markets. After centuries of unsustainable use, as well as rapidly escalating forest destruction since 1950, less than 60,000 km2 (23,000 sq mi) or 10% of Madagascar's land area remains forested. Only 17,000 km2 (6,600 sq mi) or 3% of the island's land area is protected and due to dire economic conditions and political instability, most of the protected areas are ineffectively managed and defended. Some protected areas were set aside because they were naturally protected by their remote, isolated location, often on steep cliffs. Other areas, such as the dry forests and spiny forests of the west and south, receive little protection and are in serious danger of being destroyed.
Some species may be in risk of extinction even without complete deforestation, such as ruffed lemurs, which are very sensitive to habitat disturbance. If large fruit trees are removed, the forest may sustain fewer individuals of a species and their reproductive success may be affected for years. Small populations may be able to persist in isolated forest fragments for 20 to 40 years due to long generation times, but in the long term, such populations may not be viable. Small, isolated populations also risk extirpation by natural disasters and disease outbreaks (epizootics). Two diseases that are lethal to lemurs and could severely impact isolated lemur populations are toxoplasmosis, which is spread by feral cats, and the herpes simplex virus carried by humans.
Climate change and weather-related natural disasters also threaten lemur survival. For the last 1000 years, western and highland regions have been growing significantly drier, but in the past few decades, severe drought has become much more frequent. There are indications that deforestation and forest fragmentation are accelerating this gradual desiccation. The effects of drought are even felt in the rainforests. As annual rainfall decreases, the larger trees that make up the high canopy suffer increased mortality, failure to fruit, and decreased production of new leaves, which folivorous lemurs prefer. Cyclones can defoliate an area, knock down canopy trees, and create landslides and flooding. This can leave lemur populations without fruit or leaves until the following spring, requiring them to subsist on crisis foods, such as epiphytes.
Lemurs are hunted for food by the local Malagasy, either for local subsistence or to supply a luxury meat market in the larger cities. Most rural Malagasy do not understand what "endangered" means, nor do they know that hunting lemurs is illegal or that lemurs are found only in Madagascar. Many Malagasy have taboo, or fady, about hunting and eating lemurs, but this does not prevent hunting in many regions. Even though hunting has been a threat to lemur populations in the past, it has recently become a more serious threat as socioeconomic conditions deteriorate. Economic hardships have caused people to move around the country in search of employment, leading local traditions to break down. Drought and famine can also relax the fady that protect lemurs. Larger species, such as sifakas and ruffed lemurs, are common targets, but smaller species are also hunted or accidentally caught in snares intended for larger prey. Experienced, organized hunting parties using firearms, slings and blowguns can kill as many as eight to twenty lemurs in one trip. Organized hunting parties and lemur traps can be found in both non-protected areas and remote corners of protected areas. National parks and other protected areas are not adequately protected by law enforcement agencies. Often, there are too few park rangers to cover a large area, and sometimes terrain within the park is too rugged to check regularly.
Although not as significant as deforestation and hunting, some lemurs, such as crowned lemurs and other species that have successfully been kept in captivity, are occasionally kept as exotic pets by Malagasy people. Bamboo lemurs are also kept as pets, although they only survive for up to two months. Live capture for the exotic pet trade in wealthier countries is not normally considered a threat due to strict regulations controlling their export.
Conservation efforts
Lemurs have drawn much attention to Madagascar and its endangered species. In this capacity, they act as flagship species, the most notable of which is the ring-tailed lemur, which is considered an icon of the country. The presence of lemurs in national parks helps drive ecotourism, which especially helps local communities living in the vicinity of the national parks, since it offers employment opportunities and the community receives half of the park entrance fees. In the case of Ranomafana National Park, job opportunities and other revenue from long-term research can rival that of ecotourism.
Starting in 1927, the Malagasy government has declared all lemurs as "protected" by establishing protected areas that are now classified under three categories: National Parks (Parcs Nationaux), Strict Nature Reserves (Réserves Naturelles Intégrales), and Special Reserves (Réserves Spéciales). There are currently 18 national parks, 5 strict nature reserves, and 22 special reserves, as well as several other small private reserves, such as Berenty Reserve and Sainte Luce Private Reserve, both near Fort Dauphin. All protected areas, excluding the private reserves, comprise approximately 3% of the land surface of Madagascar and are managed by Madagascar National Parks, formerly known as l'Association Nationale pour la Gestion des Aires Protégées (ANGAP), as well as other non-governmental organizations (NGOs), including Conservation International (CI), the Wildlife Conservation Society (WCS), and the World Wide Fund for Nature (WWF). Most lemur species are covered by this network of protected areas, and a few species can be found in multiple parks or reserves.
Conservation is also facilitated by the Madagascar Fauna Group (MFG), an association of nearly 40 zoos and related organizations, including the Duke Lemur Center, the Durrell Wildlife Conservation Trust, and the Saint Louis Zoological Park. This international NGO supports Madagascar's Parc Ivoloina, helps protect Betampona Reserve and other protected areas, and promotes field research, breeding programs, conservation planning, and education in zoos. One of their major projects involved the release of captive black-and-white ruffed lemurs, designed to help restock the dwindling population within Betampona Reserve.
Habitat corridors are needed for linking these protected areas so that small populations are not isolated. In in Durban, South Africa, Madagascar's former president Marc Ravalomanana promised to triple the size of the island's protected areas in five years. This became known as the "Durban Vision". In June 2007, the World Heritage Committee included a sizable portion of Madagascar's eastern rainforests as a new UNESCO World Heritage Site.
Debt relief may help Madagascar protect its biodiversity.
Captive lemur populations are maintained locally and outside of Madagascar in varied zoological conservatories and research centers, although the diversity of species is limited. Sikafas, for instance, do not survive well in captivity, so few facilities have them. The largest captive lemur population can be found at the Duke Lemur Center (DLC) in North Carolina, whose mission includes non-invasive research, conservation (e.g. captive breeding), and public education. Another large lemur colony is located in the Myakka City Lemur Reserve in Florida, run by the Lemur Conservation Foundation (LCF), which also hosts lemur research. In Madagascar, Lemurs' Park is a free-range, private facility southwest of Antananarivo that exhibits lemurs for the public while also rehabilitating captive-born lemurs for reintroduction into the wild.
In Malagasy culture
In Malagasy culture, lemurs, and animals in general, have souls (ambiroa) which can get revenge if mocked while alive or if killed in a cruel fashion. Because of this, lemurs, like many other elements of daily life, have been a source of taboos, known locally as fady, which can be based around stories with four basic principles. A village or region may believe that a certain type of lemur may be the ancestor of the clan. They may also believe that a lemur's spirit may get revenge. Alternatively, the animal may appear as a benefactor. Lemurs are also thought to impart their qualities, good or bad, onto human babies. In general, fady extend beyond a sense of the forbidden, but can include events that bring bad luck.
One example of lemur fady told around 1970 comes from Ambatofinandrahana in the Fianarantsoa Province. According to the account, a man brought a lemur home in a trap, but alive. His children wanted to keep the lemur as a pet, but when the father told them it was not a domestic animal, the children asked to kill it. After the children tortured the lemur, it eventually died and was eaten. A short time later, all the children died of illness. As a result, the father declared that anyone who tortures lemurs for fun shall "be destroyed and have no descendants."
Fady can not only help protect lemurs and their forests under stable socioeconomic situations, but they can also lead to discrimination and persecution if a lemur is known to bring bad fortune, for instance, if it walks through town. In other ways, fady does not protect all lemurs equally. For example, although the hunting and eating of certain species may be taboo, other species may not share that same protection and are therefore targeted instead. Fady can vary from village to village within the same region. If people move to a new village or region, their fady may not apply to the lemur species that are locally present, making them available for consumption. Fady restrictions on lemur meat can be relaxed in times of famine and drought.
The aye-aye is almost universally viewed unfavorably across Madagascar, though the tales vary from village to village and region to region. If people see an aye-aye, they may kill it and hang the corpse on a pole near a road outside of town (so others can carry the bad fortunes away) or burn their village and move. The superstitions behind aye-aye fady include beliefs that they kill and eat chickens or people, that they kill people in their sleep by cutting their aortic vein, that they embody ancestral spirits, or that they warn of illness, death, or bad luck in the family. As of 1970, the people of the Marolambo District in the Toamasina Province feared the aye-aye because they believed it had supernatural powers. Because of this, no one was allowed to mock, kill, or eat one.
There are also widespread fady about indri and sifakas. They are often protected from hunting and consumption because of their resemblance to humans and their ancestors, mostly due to their large size and upright or orthograde posture. The resemblance is even stronger for indri, which lack the long tail of most living lemurs. Known locally as babakoto ("Ancestor of Man"), the indri is sometimes seen as the progenitor of the family or clan. There are also stories of an indri that helped a human down from a tree, so they are seen as benefactors. Other lemur fady include the belief that a wife will have ugly children if her husband kills a woolly lemur, or that if a pregnant woman eats a dwarf lemur, her baby will get its beautiful, round eyes.
In popular culture
Lemurs have also become popular in Western culture in recent years. The DreamWorks Animation franchise Madagascar features the characters King Julien, Maurice and Mort and was seen by an estimated 100 million people in theaters and 200–300 million people on DVD worldwide. Prior to this film, Zoboomafoo, a Public Broadcasting Service (PBS) children's television series from 1999 to 2001, helped to popularize sifakas by featuring a live Coquerel's sifaka from the Duke Lemur Center as well as a puppet. A twenty-episode series called Lemur Kingdom (in the United States) or Lemur Street (in the United Kingdom and Canada) aired in 2008 on Animal Planet. It combined the typical animal documentary with dramatic narration to tell the story of two groups of ring-tailed lemurs at Berenty Private Reserve.
| Biology and health sciences | Primates | null |
478091 | https://en.wikipedia.org/wiki/Von%20Neumann%20architecture | Von Neumann architecture | The von Neumann architecture—also known as the von Neumann model or Princeton architecture—is a computer architecture based on the First Draft of a Report on the EDVAC, written by John von Neumann in 1945, describing designs discussed with John Mauchly and J. Presper Eckert at the University of Pennsylvania's Moore School of Electrical Engineering. The document describes a design architecture for an electronic digital computer with these components:
A processing unit with both an arithmetic logic unit and processor registers
A control unit that includes an instruction register and a program counter
Memory that stores data and instructions
External mass storage
Input and output mechanisms
The attribution of the invention of the architecture to von Neumann is controversial, not least because Eckert and Mauchly had done a lot of the required design work and claim to have had the idea for stored programs long before discussing the ideas with von Neumann and Herman Goldstine.
The term "von Neumann architecture" has evolved to refer to any stored-program computer in which an instruction fetch and a data operation cannot occur at the same time (since they share a common bus). This is referred to as the von Neumann bottleneck, which often limits the performance of the corresponding system.
The von Neumann architecture is simpler than the Harvard architecture (which has one dedicated set of address and data buses for reading and writing to memory and another set of address and data buses to fetch instructions).
A stored-program computer uses the same underlying mechanism to encode both program instructions and data as opposed to designs which use a mechanism such as discrete plugboard wiring or fixed control circuitry for instruction implementation. Stored-program computers were an advancement over the manually reconfigured or fixed function computers of the 1940s, such as the Colossus and the ENIAC. These were programmed by setting switches and inserting patch cables to route data and control signals between various functional units.
The vast majority of modern computers use the same hardware mechanism to encode and store both data and program instructions, but have caches between the CPU and memory, and, for the caches closest to the CPU, have separate caches for instructions and data, so that most instruction and data fetches use separate buses (split-cache architecture).
History
The earliest computing machines had fixed programs. Some very simple computers still use this design, either for simplicity or training purposes. For example, a desk calculator (in principle) is a fixed program computer. It can do basic mathematics, but it cannot run a word processor or games. Changing the program of a fixed-program machine requires rewiring, restructuring, or redesigning the machine. The earliest computers were not so much "programmed" as "designed" for a particular task. "Reprogramming"—when possible at all—was a laborious process that started with flowcharts and paper notes, followed by detailed engineering designs, and then the often-arduous process of physically rewiring and rebuilding the machine. It could take three weeks to set up and debug a program on ENIAC.
With the proposal of the stored-program computer, this changed. A stored-program computer includes, by design, an instruction set, and can store in memory a set of instructions (a program) that details the computation.
A stored-program design also allows for self-modifying code. One early motivation for such a facility was the need for a program to increment or otherwise modify the address portion of instructions, which operators had to do manually in early designs. This became less important when index registers and indirect addressing became usual features of machine architecture. Another use was to embed frequently used data in the instruction stream using immediate addressing.
Capabilities
On a large scale, the ability to treat instructions as data is what makes assemblers, compilers, linkers, loaders, and other automated programming tools possible. It makes "programs that write programs" possible. This has made a sophisticated self-hosting computing ecosystem flourish around von Neumann architecture machines.
Some high-level languages leverage the von Neumann architecture by providing an abstract, machine-independent way to manipulate executable code at runtime (e.g., LISP), or by using runtime information to tune just-in-time compilation (e.g. languages hosted on the Java virtual machine, or languages embedded in web browsers).
On a smaller scale, some repetitive operations such as BITBLT or pixel and vertex shaders can be accelerated on general purpose processors with just-in-time compilation techniques. This is one use of self-modifying code that has remained popular.
Development of the stored-program concept
The mathematician Alan Turing, who had been alerted to a problem of mathematical logic by the lectures of Max Newman at the University of Cambridge, wrote a paper in 1936 entitled On Computable Numbers, with an Application to the Entscheidungsproblem, which was published in the Proceedings of the London Mathematical Society. In it he described a hypothetical machine he called a universal computing machine, now known as the "Universal Turing machine". The hypothetical machine had an infinite store (memory in today's terminology) that contained both instructions and data. John von Neumann became acquainted with Turing while he was a visiting professor at Cambridge in 1935, and also during Turing's PhD year at the Institute for Advanced Study in Princeton, New Jersey during 1936–1937. Whether he knew of Turing's paper of 1936 at that time is not clear.
In 1936, Konrad Zuse also anticipated, in two patent applications, that machine instructions could be stored in the same storage used for data.
Independently, J. Presper Eckert and John Mauchly, who were developing the ENIAC at the Moore School of Electrical Engineering of the University of Pennsylvania, wrote about the stored-program concept in December 1943.
In planning a new machine, EDVAC, Eckert wrote in January 1944 that they would store data and programs in a new addressable memory device, a mercury metal delay-line memory. This was the first time the construction of a practical stored-program machine was proposed. At that time, he and Mauchly were not aware of Turing's work.
Von Neumann was involved in the Manhattan Project at the Los Alamos National Laboratory. It required huge amounts of calculation, and thus drew him to the ENIAC project, during the summer of 1944. There he joined the ongoing discussions on the design of this stored-program computer, the EDVAC. As part of that group, he wrote up a description titled First Draft of a Report on the EDVAC based on the work of Eckert and Mauchly. It was unfinished when his colleague Herman Goldstine circulated it, and bore only von Neumann's name (to the consternation of Eckert and Mauchly). The paper was read by dozens of von Neumann's colleagues in America and Europe, and influenced the next round of computer designs.
Jack Copeland considers that it is "historically inappropriate to refer to electronic stored-program digital computers as 'von Neumann machines. His Los Alamos colleague Stan Frankel said of von Neumann's regard for Turing's ideas
At the time that the "First Draft" report was circulated, Turing was producing a report entitled Proposed Electronic Calculator. It described in engineering and programming detail, his idea of a machine he called the Automatic Computing Engine (ACE). He presented this to the executive committee of the British National Physical Laboratory on February 19, 1946. Although Turing knew from his wartime experience at Bletchley Park that what he proposed was feasible, the secrecy surrounding Colossus, that was subsequently maintained for several decades, prevented him from saying so. Various successful implementations of the ACE design were produced.
Both von Neumann's and Turing's papers described stored-program computers, but von Neumann's earlier paper achieved greater circulation and the computer architecture it outlined became known as the "von Neumann architecture". In the 1953 publication Faster than Thought: A Symposium on Digital Computing Machines (edited by B. V. Bowden), a section in the chapter on Computers in America reads as follows:
The Machine of the Institute For Advanced Study, Princeton
In 1945, Professor J. von Neumann, who was then working at the Moore School of Engineering in Philadelphia, where the E.N.I.A.C. had been built, issued on behalf of a group of his co-workers, a report on the logical design of digital computers. The report contained a detailed proposal for the design of the machine that has since become known as the E.D.V.A.C. (electronic discrete variable automatic computer). This machine has only recently been completed in America, but the von Neumann report inspired the construction of the E.D.S.A.C. (electronic delay-storage automatic calculator) in Cambridge (see p. 130).
In 1947, Burks, Goldstine and von Neumann published another report that outlined the design of another type of machine (a parallel machine this time) that would be exceedingly fast, capable perhaps of 20,000 operations per second. They pointed out that the outstanding problem in constructing such a machine was the development of suitable memory with instantaneously accessible contents. At first they suggested using a special vacuum tube—called the "Selectron"—which the Princeton Laboratories of RCA had invented. These tubes were expensive and difficult to make, so von Neumann subsequently decided to build a machine based on the Williams memory. This machine—completed in June, 1952 in Princeton—has become popularly known as the Maniac. The design of this machine inspired at least half a dozen machines now being built in America, all known affectionately as "Johniacs".
In the same book, the first two paragraphs of a chapter on ACE read as follows:
Automatic Computation at the National Physical Laboratory
One of the most modern digital computers which embodies developments and improvements in the technique of automatic electronic computing was recently demonstrated at the National Physical Laboratory, Teddington, where it has been designed and built by a small team of mathematicians and electronics research engineers on the staff of the Laboratory, assisted by a number of production engineers from the English Electric Company, Limited. The equipment so far erected at the Laboratory is only the pilot model of a much larger installation which will be known as the Automatic Computing Engine, but although comparatively small in bulk and containing only about 800 thermionic valves, as can be judged from Plates XII, XIII and XIV, it is an extremely rapid and versatile calculating machine.
The basic concepts and abstract principles of computation by a machine were formulated by Dr. A. M. Turing, F.R.S., in a paper1. read before the London Mathematical Society in 1936, but work on such machines in Britain was delayed by the war. In 1945, however, an examination of the problems was made at the National Physical Laboratory by Mr. J. R. Womersley, then superintendent of the Mathematics Division of the Laboratory. He was joined by Dr. Turing and a small staff of specialists, and, by 1947, the preliminary planning was sufficiently advanced to warrant the establishment of the special group already mentioned. In April, 1948, the latter became the Electronics Section of the Laboratory, under the charge of Mr. F. M. Colebrook.
Early von Neumann-architecture computers
The First Draft described a design that was used by many universities and corporations to construct their computers. Among these various computers, only ILLIAC and ORDVAC had compatible instruction sets.
ARC2 (Birkbeck, University of London) officially came online on May 12, 1948.
Manchester Baby (Victoria University of Manchester, England) made its first successful run of a stored program on June 21, 1948.
EDSAC (University of Cambridge, England) was the first practical stored-program electronic computer (May 1949)
Manchester Mark 1 (University of Manchester, England) Developed from the Baby (June 1949)
CSIRAC (Council for Scientific and Industrial Research) Australia (November 1949)
MESM at the Kiev Institute of Electrotechnology in Kiev, Ukrainian SSR (November 1950)
EDVAC (Ballistic Research Laboratory, Computing Laboratory at Aberdeen Proving Ground 1951)
IAS machine at Institute for Advanced Study (1951)
ORDVAC (University of Illinois) at Aberdeen Proving Ground, Maryland (completed November 1951)
MANIAC I at Los Alamos Scientific Laboratory (March 1952)
ILLIAC at the University of Illinois, (September 1952)
BESM-1 in Moscow (1952)
AVIDAC at Argonne National Laboratory (1953)
ORACLE at Oak Ridge National Laboratory (June 1953)
BESK in Stockholm (1953)
JOHNNIAC at RAND Corporation (January 1954)
DASK in Denmark (1955)
WEIZAC at the Weizmann Institute of Science in Rehovot, Israel (1955)
PERM in Munich (1956)
SILLIAC in Sydney (1956)
Early stored-program computers
The date information in the following chronology is difficult to put into proper order. Some dates are for first running a test program, some dates are the first time the computer was demonstrated or completed, and some dates are for the first delivery or installation.
The IBM SSEC had the ability to treat instructions as data, and was publicly demonstrated on January 27, 1948. This ability was claimed in a US patent. However it was partially electromechanical, not fully electronic. In practice, instructions were read from paper tape due to its limited memory.
The ARC2 developed by Andrew Booth and Kathleen Booth at Birkbeck, University of London officially came online on May 12, 1948. It featured the first rotating drum storage device.
The Manchester Baby was the first fully electronic computer to run a stored program. It ran a factoring program for 52 minutes on June 21, 1948, after running a simple division program and a program to show that two numbers were relatively prime.
The ENIAC was modified to run as a primitive read-only stored-program computer (using the Function Tables for program ROM) and was demonstrated as such on September 16, 1948, running a program by Adele Goldstine for von Neumann.
The BINAC ran some test programs in February, March, and April 1949, although was not completed until September 1949.
The Manchester Mark 1 developed from the Baby project. An intermediate version of the Mark 1 was available to run programs in April 1949, but was not completed until October 1949.
The EDSAC ran its first program on May 6, 1949.
The EDVAC was delivered in August 1949, but it had problems that kept it from being put into regular operation until 1951.
The CSIR Mk I ran its first program in November 1949.
The SEAC was demonstrated in April 1950.
The Pilot ACE ran its first program on May 10, 1950, and was demonstrated in December 1950.
The SWAC was completed in July 1950.
The Whirlwind was completed in December 1950 and was in actual use in April 1951.
The first ERA Atlas (later the commercial ERA 1101/UNIVAC 1101) was installed in December 1950.
Evolution
Through the decades of the 1960s and 1970s computers generally became both smaller and faster, which led to evolutions in their architecture. For example, memory-mapped I/O lets input and output devices be treated the same as memory. A single system bus could be used to provide a modular system with lower cost. This is sometimes called a "streamlining" of the architecture.
In subsequent decades, simple microcontrollers would sometimes omit features of the model to lower cost and size.
Larger computers added features for higher performance.
Design limitations
von Neumann bottleneck
The use of the same bus to fetch instructions and data leads to the von Neumann bottleneck, the limited throughput (data transfer rate) between the central processing unit (CPU) and memory compared to the amount of memory. Because the single bus can only access one of the two classes of memory at a time, throughput is lower than the rate at which the CPU can work. This seriously limits the effective processing speed when the CPU is required to perform minimal processing on large amounts of data. The CPU is continually forced to wait for needed data to move to or from memory. Since CPU speed and memory size have increased much faster than the throughput between them, the bottleneck has become more of a problem, a problem whose severity increases with every new generation of CPU.
The von Neumann bottleneck was described by John Backus in his 1977 ACM Turing Award lecture. According to Backus:
Surely there must be a less primitive way of making big changes in the store than by pushing vast numbers of words back and forth through the von Neumann bottleneck. Not only is this tube a literal bottleneck for the data traffic of a problem, but, more importantly, it is an intellectual bottleneck that has kept us tied to word-at-a-time thinking instead of encouraging us to think in terms of the larger conceptual units of the task at hand. Thus programming is basically planning and detailing the enormous traffic of words through the von Neumann bottleneck, and much of that traffic concerns not significant data itself, but where to find it.
Mitigations
There are several known methods for mitigating the Von Neumann performance bottleneck. For example, the following all can improve performance:
providing a cache between the CPU and the main memory;
providing separate caches or separate access paths for data and instructions (the so-called Modified Harvard architecture);
using branch predictor algorithms and logic;
providing a limited CPU stack or other on-chip scratchpad memory to reduce memory access;
implementing the CPU and the memory hierarchy as a system on chip, providing greater locality of reference and thus reducing latency and increasing throughput between processor registers and main memory.
The problem can also be sidestepped somewhat by using parallel computing, using for example the non-uniform memory access (NUMA) architecture—this approach is commonly employed by supercomputers. It is less clear whether the intellectual bottleneck that Backus criticized has changed much since 1977. Backus's proposed solution has not had a major influence. Modern functional programming and object-oriented programming are much less geared towards "pushing vast numbers of words back and forth" than earlier languages like FORTRAN were, but internally, that is still what computers spend much of their time doing, even highly parallel supercomputers.
As of 1996, a database benchmark study found that three out of four CPU cycles were spent waiting for memory. Researchers expect that increasing the number of simultaneous instruction streams with multithreading or single-chip multiprocessing will make this bottleneck even worse. In the context of multi-core processors, additional overhead is required to maintain cache coherence between processors and threads.
Self-modifying code
Aside from the von Neumann bottleneck, program modifications can be quite harmful, either by accident or design. In some simple stored-program computer designs, a malfunctioning program can damage itself, other programs, or the operating system, possibly leading to a computer crash. However, this problem also applies to conventional programs that lack bounds checking. Memory protection and various access controls generally safeguard against both accidental and malicious program changes.
| Technology | Computer architecture concepts | null |
478195 | https://en.wikipedia.org/wiki/Normal%20mode | Normal mode | A normal mode of a dynamical system is a pattern of motion in which all parts of the system move sinusoidally with the same frequency and with a fixed phase relation. The free motion described by the normal modes takes place at fixed frequencies. These fixed frequencies of the normal modes of a system are known as its natural frequencies or resonant frequencies. A physical object, such as a building, bridge, or molecule, has a set of normal modes and their natural frequencies that depend on its structure, materials and boundary conditions.
The most general motion of a linear system is a superposition of its normal modes. The modes are normal in the sense that they can move independently, that is to say that an excitation of one mode will never cause motion of a different mode. In mathematical terms, normal modes are orthogonal to each other.
General definitions
Mode
In the wave theory of physics and engineering, a mode in a dynamical system is a standing wave state of excitation, in which all the components of the system will be affected sinusoidally at a fixed frequency associated with that mode.
Because no real system can perfectly fit under the standing wave framework, the mode concept is taken as a general characterization of specific states of oscillation, thus treating the dynamic system in a linear fashion, in which linear superposition of states can be performed.
Typical examples include:
In a mechanical dynamical system, a vibrating rope is the most clear example of a mode, in which the rope is the medium, the stress on the rope is the excitation, and the displacement of the rope with respect to its static state is the modal variable.
In an acoustic dynamical system, a single sound pitch is a mode, in which the air is the medium, the sound pressure in the air is the excitation, and the displacement of the air molecules is the modal variable.
In a structural dynamical system, a high tall building oscillating under its most flexural axis is a mode, in which all the material of the building -under the proper numerical simplifications- is the medium, the seismic/wind/environmental solicitations are the excitations and the displacements are the modal variable.
In an electrical dynamical system, a resonant cavity made of thin metal walls, enclosing a hollow space, for a particle accelerator is a pure standing wave system, and thus an example of a mode, in which the hollow space of the cavity is the medium, the RF source (a Klystron or another RF source) is the excitation and the electromagnetic field is the modal variable.
When relating to music, normal modes of vibrating instruments (strings, air pipes, drums, etc.) are called "overtones".
The concept of normal modes also finds application in other dynamical systems, such as optics, quantum mechanics, atmospheric dynamics and molecular dynamics.
Most dynamical systems can be excited in several modes, possibly simultaneously. Each mode is characterized by one or several frequencies, according to the modal variable field. For example, a vibrating rope in 2D space is defined by a single-frequency (1D axial displacement), but a vibrating rope in 3D space is defined by two frequencies (2D axial displacement).
For a given amplitude on the modal variable, each mode will store a specific amount of energy because of the sinusoidal excitation.
The normal or dominant mode of a system with multiple modes will be the mode storing the minimum amount of energy for a given amplitude of the modal variable, or, equivalently, for a given stored amount of energy, the dominant mode will be the mode imposing the maximum amplitude of the modal variable.
Mode numbers
A mode of vibration is characterized by a modal frequency and a mode shape. It is numbered according to the number of half waves in the vibration. For example, if a vibrating beam with both ends pinned displayed a mode shape of half of a sine wave (one peak on the vibrating beam) it would be vibrating in mode 1. If it had a full sine wave (one peak and one trough) it would be vibrating in mode 2.
In a system with two or more dimensions, such as the pictured disk, each dimension is given a mode number. Using polar coordinates, we have a radial coordinate and an angular coordinate. If one measured from the center outward along the radial coordinate one would encounter a full wave, so the mode number in the radial direction is 2. The other direction is trickier, because only half of the disk is considered due to the anti-symmetric (also called skew-symmetry) nature of a disk's vibration in the angular direction. Thus, measuring 180° along the angular direction you would encounter a half wave, so the mode number in the angular direction is 1. So the mode number of the system is 2–1 or 1–2, depending on which coordinate is considered the "first" and which is considered the "second" coordinate (so it is important to always indicate which mode number matches with each coordinate direction).
In linear systems each mode is entirely independent of all other modes. In general all modes have different frequencies (with lower modes having lower frequencies) and different mode shapes.
Nodes
In a one-dimensional system at a given mode the vibration will have nodes, or places where the displacement is always zero. These nodes correspond to points in the mode shape where the mode shape is zero. Since the vibration of a system is given by the mode shape multiplied by a time function, the displacement of the node points remain zero at all times.
When expanded to a two dimensional system, these nodes become lines where the displacement is always zero. If you watch the animation above you will see two circles (one about halfway between the edge and center, and the other on the edge itself) and a straight line bisecting the disk, where the displacement is close to zero. In an idealized system these lines equal zero exactly, as shown to the right.
In mechanical systems
In the analysis of conservative systems with small displacements from equilibrium, important in acoustics, molecular spectra, and electrical circuits, the system can be transformed to new coordinates called normal coordinates. Each normal coordinate corresponds to a single vibrational frequency of the system and the corresponding motion of the system is called the normal mode of vibration.
Coupled oscillators
Consider two equal bodies (not affected by gravity), each of mass , attached to three springs, each with spring constant . They are attached in the following manner, forming a system that is physically symmetric:
where the edge points are fixed and cannot move. Let denote the horizontal displacement of the left mass, and denote the displacement of the right mass.
Denoting acceleration (the second derivative of with respect to time) as the equations of motion are:
Since we expect oscillatory motion of a normal mode (where is the same for both masses), we try:
Substituting these into the equations of motion gives us:
Omitting the exponential factor (because it is common to all terms) and simplifying yields:
And in matrix representation:
If the matrix on the left is invertible, the unique solution is the trivial solution . The non trivial solutions are to be found for those values of whereby the matrix on the left is singular; i.e. is not invertible. It follows that the determinant of the matrix must be equal to 0, so:
Solving for , the two positive solutions are:
Substituting into the matrix and solving for , yields . Substituting results in . (These vectors are eigenvectors, and the frequencies are eigenvalues.)
The first normal mode is:
Which corresponds to both masses moving in the same direction at the same time. This mode is called antisymmetric.
The second normal mode is:
This corresponds to the masses moving in the opposite directions, while the center of mass remains stationary. This mode is called symmetric.
The general solution is a superposition of the normal modes where , , , and are determined by the initial conditions of the problem.
The process demonstrated here can be generalized and formulated using the formalism of Lagrangian mechanics or Hamiltonian mechanics.
Standing waves
A standing wave is a continuous form of normal mode. In a standing wave, all the space elements (i.e. coordinates) are oscillating in the same frequency and in phase (reaching the equilibrium point together), but each has a different amplitude.
The general form of a standing wave is:
where represents the dependence of amplitude on location and the cosine/sine are the oscillations in time.
Physically, standing waves are formed by the interference (superposition) of waves and their reflections (although one may also say the opposite; that a moving wave is a superposition of standing waves). The geometric shape of the medium determines what would be the interference pattern, thus determines the form of the standing wave. This space-dependence is called a normal mode.
Usually, for problems with continuous dependence on there is no single or finite number of normal modes, but there are infinitely many normal modes. If the problem is bounded (i.e. it is defined on a finite section of space) there are countably many normal modes (usually numbered ). If the problem is not bounded, there is a continuous spectrum of normal modes.
Elastic solids
In any solid at any temperature, the primary particles (e.g. atoms or molecules) are not stationary, but rather vibrate about mean positions. In insulators the capacity of the solid to store thermal energy is due almost entirely to these vibrations. Many physical properties of the solid (e.g. modulus of elasticity) can be predicted given knowledge of the frequencies with which the particles vibrate. The simplest assumption (by Einstein) is that all the particles oscillate about their mean positions with the same natural frequency . This is equivalent to the assumption that all atoms vibrate independently with a frequency . Einstein also assumed that the allowed energy states of these oscillations are harmonics, or integral multiples of . The spectrum of waveforms can be described mathematically using a Fourier series of sinusoidal density fluctuations (or thermal phonons).
Debye subsequently recognized that each oscillator is intimately coupled to its neighboring oscillators at all times. Thus, by replacing Einstein's identical uncoupled oscillators with the same number of coupled oscillators, Debye correlated the elastic vibrations of a one-dimensional solid with the number of mathematically special modes of vibration of a stretched string (see figure). The pure tone of lowest pitch or frequency is referred to as the fundamental and the multiples of that frequency are called its harmonic overtones. He assigned to one of the oscillators the frequency of the fundamental vibration of the whole block of solid. He assigned to the remaining oscillators the frequencies of the harmonics of that fundamental, with the highest of all these frequencies being limited by the motion of the smallest primary unit.
The normal modes of vibration of a crystal are in general superpositions of many overtones, each with an appropriate amplitude and phase. Longer wavelength (low frequency) phonons are exactly those acoustical vibrations which are considered in the theory of sound. Both longitudinal and transverse waves can be propagated through a solid, while, in general, only longitudinal waves are supported by fluids.
In the longitudinal mode, the displacement of particles from their positions of equilibrium coincides with the propagation direction of the wave. Mechanical longitudinal waves have been also referred to as . For transverse modes, individual particles move perpendicular to the propagation of the wave.
According to quantum theory, the mean energy of a normal vibrational mode of a crystalline solid with characteristic frequency is:
The term represents the "zero-point energy", or the energy which an oscillator will have at absolute zero. tends to the classic value at high temperatures
By knowing the thermodynamic formula,
the entropy per normal mode is:
The free energy is:
which, for , tends to:
In order to calculate the internal energy and the specific heat, we must know the number of normal vibrational modes a frequency between the values and . Allow this number to be . Since the total number of normal modes is , the function is given by:
The integration is performed over all frequencies of the crystal. Then the internal energy will be given by:
In quantum mechanics
Bound states
in quantum mechanics are analogous to modes. The waves in quantum systems are oscillations in probability amplitude rather than material displacement. The frequency of oscillation, , relates to the mode energy by where is the Planck constant. Thus a system like an atom consists of a linear combination of modes of definite energy. These energies are characteristic of the particular atom. The (complex) square of the probability amplitude at a point in space gives the probability of measuring an electron at that location. The spatial distribution of this probability is characteristic of the atom.
In seismology
Normal modes are generated in the Earth from long wavelength seismic waves from large earthquakes interfering to form standing waves.
For an elastic, isotropic, homogeneous sphere, spheroidal, toroidal and radial (or breathing) modes arise. Spheroidal modes only involve P and SV waves (like Rayleigh waves) and depend on overtone number and angular order but have degeneracy of azimuthal order . Increasing concentrates fundamental branch closer to surface and at large this tends to Rayleigh waves. Toroidal modes only involve SH waves (like Love waves) and do not exist in fluid outer core. Radial modes are just a subset of spheroidal modes with . The degeneracy does not exist on Earth as it is broken by rotation, ellipticity and 3D heterogeneous velocity and density structure.
It may be assumed that each mode can be isolated, the self-coupling approximation, or that many modes close in frequency resonate, the cross-coupling approximation. Self-coupling will solely change the phase velocity and not the number of waves around a great circle, resulting in a stretching or shrinking of standing wave pattern. Modal cross-coupling occurs due to the rotation of the Earth, from aspherical elastic structure, or due to Earth's ellipticity and leads to a mixing of fundamental spheroidal and toroidal modes.
| Physical sciences | Basics_4 | Physics |
478224 | https://en.wikipedia.org/wiki/Orcus%20%28dwarf%20planet%29 | Orcus (dwarf planet) | Orcus (minor-planet designation: 90482 Orcus) is a dwarf planet located in the Kuiper belt, with one large moon, Vanth. It has an estimated diameter of , comparable to the Inner Solar System dwarf planet Ceres. The surface of Orcus is relatively bright with albedo reaching 23 percent, neutral in color, and rich in water ice. The ice is predominantly in crystalline form, which may be related to past cryovolcanic activity. Other compounds like methane or ammonia may also be present on its surface. Orcus was discovered by American astronomers Michael Brown, Chad Trujillo, and David Rabinowitz on 17 February 2004.
Orcus is a plutino, a trans-Neptunian object that is locked in a 2:3 orbital resonance with the ice giant Neptune, making two revolutions around the Sun to every three of Neptune's. This is much like Pluto, except that the phase of Orcus's orbit is opposite to Pluto's: Orcus is at aphelion (most recently in 2019) around when Pluto is at perihelion (most recently in 1989) and vice versa. Orcus is the second-largest known plutino, after Pluto itself. The perihelion of Orcus's orbit is around 120° from that of Pluto, while the eccentricities and inclinations are similar. Because of these similarities and contrasts, along with its large moon Vanth that can be compared to Pluto's large moon Charon, Orcus has been dubbed the "." This was a major consideration in selecting its name, as the deity Orcus was the Roman/Etruscan equivalent of the Roman/Greek Pluto.
History
Discovery
Orcus was discovered on 17 February 2004, by American astronomers Michael Brown of Caltech, Chad Trujillo of the Gemini Observatory, and David Rabinowitz of Yale University. Precovery images taken by the Palomar Observatory as early as 8 November 1951 were later obtained from the Digitized Sky Survey.
Name and symbol
The minor planet Orcus was named after one of the Roman gods of the underworld, Orcus. While Pluto (of Greek origin) was the ruler of the underworld, Orcus (of Etruscan origin) was a punisher of the condemned. The name was published by the Minor Planet Center on 26 November 2004 (). Under the guidelines of the International Astronomical Union's (IAU) naming conventions, objects with a similar size and orbit to that of Pluto are named after underworld deities. Accordingly, the discoverers suggested naming the object after Orcus, the Etruscan god of the underworld and punisher of broken oaths. The name was also a private reference to the homonymous Orcas Island, where Brown's wife had lived as a child and that they visit frequently.
On 30 March 2005, Orcus's moon, Vanth, was named after a winged female entity, Vanth, of the Etruscan underworld. She could be present at the moment of death, and frequently acted as a psychopomp, a guide of the deceased to the underworld.
The usage of planetary symbols is no longer recommended in astronomy, so Orcus never received a symbol in the astronomical literature. A symbol , used mostly among astrologers, is included in Unicode as U+1F77F. The symbol was designed by Denis Moskowitz, a software engineer in Massachusetts; it is an OR monogram, designed to resemble both a skull and an orca's gape. There is a rarer symbol , an inverted astrological Pluto symbol, reflecting Orcus as the anti-Pluto: it was designed by Melanie Reinhart.
Orbit and rotation
Orcus is in a 2:3 orbital resonance with Neptune, having an orbital period of 245 years, and is classified as a plutino. Its orbit is moderately inclined at 20.6° to the ecliptic. Orcus's orbit is similar to Pluto's (both have perihelia above the ecliptic), but is oriented differently. Although at one point its orbit approaches that of Neptune, the resonance between the two bodies means that Orcus itself is always a great distance away from Neptune (there is always an angular separation of over 60° between them). Over a 14,000-year period, Orcus stays more than 18 AU from Neptune. Because their mutual resonance with Neptune constrains Orcus and Pluto to remain in opposite phases of their otherwise very similar motions, Orcus is sometimes described as the "anti-Pluto". Orcus last reached its aphelion (farthest distance from the Sun) in 2019 and will come to perihelion (closest distance to the Sun) around 10 January 2143. Simulations by the Deep Ecliptic Survey show that over the next 10 million years Orcus may acquire a perihelion distance (qmin) as small as 27.8 AU.
The rotation period of Orcus is uncertain, as different photometric surveys have produced different results. Some show low amplitude variations with periods ranging from 7 to 21 hours, whereas others show no variability. The rotational axis of Orcus probably coincides with the orbital axis of its moon, Vanth. This means that Orcus is currently viewed pole-on, which could explain the near absence of any rotational modulation of its brightness. Astronomer José Luis Ortiz and colleagues have derived a possible rotation period of about 10.5 hours, assuming that Orcus is not tidally locked with Vanth. If, however, the primary is tidally locked with the satellite, the rotational period would coincide with the 9.7-day orbital period of Vanth.
Physical characteristics
Size and magnitude
The absolute magnitude of Orcus is approximately 2.3. The detection of Orcus by the Spitzer Space Telescope in the far infrared and by Herschel Space Telescope in submillimeter estimates its diameter at , with an uncertainty of . Orcus appears to have an albedo of about 21–25%, which may be typical of trans-Neptunian objects approaching the diameter range. The magnitude and size estimates were made under the assumption that Orcus is a singular object. The presence of a relatively large satellite, Vanth, may change them considerably. The absolute magnitude of Vanth is estimated at 4.88, which means that it is about 1/11 as bright as Orcus itself. The ALMA submillimeter measurements taken in 2016 showed that Vanth has a relatively large size of with an albedo of about 8 percent while Orcus's has a slightly smaller size of . Using a stellar occultation by Vanth in 2017, Vanth's diameter has been determined to be , with an uncertainty of . Michael Brown's website lists Orcus as a dwarf planet with "near certainty", Tancredi concludes that it is one, and is massive enough to be considered one under the 2006 draft proposal of the IAU, but the IAU has not formally recognized it as such.
Mass and density
Orcus and Vanth are known to constitute a binary system. The mass of the system has been estimated to be , approximately equal to that of the Saturnian moon Tethys (). The mass of the Orcus system is about 3.8 percent that of , the most massive known dwarf planet ().
The ratio of the mass of Vanth to that of Orcus was measured astrometrically with the ALMA submillimeter telescope and is with Vanth containing of the total system mass. This also means that the densities of both bodies are about the same at ~.
Spectra and surface
The first spectroscopic observations in 2004 showed that the visible spectrum of Orcus is flat (neutral in color) and featureless, whereas in the near-infrared there were moderately strong water absorption bands at 1.5 and 2.0 μm. The neutral visible spectrum and strong water absorption bands of Orcus showed that Orcus appeared different from other trans-Neptunian objects, which typically have a red visible spectrum and often featureless infrared spectra. Further infrared observations in 2004 by the European Southern Observatory and the Gemini telescope gave results consistent with mixtures of water ice and carbonaceous compounds, such as tholins. The water and methane ices can cover no more than 50 percent and 30 percent of the surface, respectively. This means the proportion of ice on the surface is less than on Charon, but similar to that on Triton.
Later in 2008–2010 new infrared spectroscopic observations with a higher signal-to-noise ratio revealed additional spectral features. Among them is a deep water ice absorption band at 1.65 μm, which is evidence of the crystalline water ice on the surface of Orcus, and a new absorption band at 2.22 μm. The origin of the latter feature is not completely clear. It can be caused either by ammonia/ammonium dissolved in the water ice or by methane/ethane ices. The radiative transfer modeling showed that a mixture of water ice, tholins (as a darkening agent), ethane ice, and ammonium ion (NH4+) provides the best match to the spectra, whereas a combination of water ice, tholins, methane ice and ammonia hydrate gives a slightly inferior result. On the other hand, a mixture of only ammonia hydrate, tholins and water ice failed to provide a satisfactory match. As of 2010, the only reliably identified compounds on the surface of Orcus are crystalline water ice and, possibly, dark tholins. A firm identification of ammonia, methane, and other hydrocarbons requires a better infrared spectra.
Orcus sits at the threshold for trans-Neptunian objects massive enough to retain volatiles such as methane on the surface. The reflectance spectrum of Orcus shows the deepest water-ice absorption bands of any Kuiper belt object that is not associated with the Haumea collisional family. The large icy satellites of Uranus have infrared spectra quite similar to that of Orcus. Among other trans-Neptunian objects, the large plutino and Pluto's moon Charon both have similar surface spectra to Orcus, with flat, featureless visible spectra and moderately strong water ice absorption bands in the near-infrared.
Cryovolcanism
Crystalline water ice on the surfaces of trans-Neptunian objects should be completely amorphized by the galactic and Solar radiation in about 10 million years. Thus the presence of crystalline water ice, and possibly ammonia ice, may indicate that a renewal mechanism was active in the past on the surface of Orcus. Ammonia so far has not been detected on any trans-Neptunian object or icy satellite of the outer planets other than Miranda. The 1.65 μm band on Orcus is broad and deep (12%), as on Charon, , , and icy satellites of the giant planets. Some calculations indicate that cryovolcanism, which is considered one of the possible renewal mechanisms, may indeed be possible for trans-Neptunian objects larger than about . Orcus may have experienced at least one such episode in the past, which turned the amorphous water ice on its surface into crystalline. The preferred type of volcanism may have been explosive aqueous volcanism driven by an explosive dissolution of methane from water–ammonia melts.
Satellite
Orcus has one known moon, Vanth (full designation ). It was discovered by Michael Brown and T.-A. Suer using discovery images taken by the Hubble Space Telescope on 13 November 2005. The discovery was announced in an IAU Circular notice published on 22 February 2007. A spatially resolved submillimeter imaging of Orcus–Vanth system in 2016 showed that Vanth has a relatively large size of , with an uncertainty of . That estimate for Vanth is in good agreement with the size of about derived from a stellar occultation in 2017. Like Charon compared to Pluto, Vanth is quite large compared to Orcus, and is one reason for characterizing Orcus as the 'anti-Pluto'. If Orcus is a dwarf planet, Vanth would be the third-largest known dwarf-planet moon, after Charon and Dysnomia. The ratio of masses of Orcus and Vanth is uncertain, possibly anywhere from 1:33 to 1:12.
| Physical sciences | Solar System | Astronomy |
478313 | https://en.wikipedia.org/wiki/Marker%20pen | Marker pen | A marker pen, fine liner, marking pen, felt-tip pen, felt pen, flow marker, sign pen (in South Korea), vivid (in New Zealand), flomaster (in East and South Slavic countries), texta (in Australia), sketch pen (in South Asia), koki (in South Africa) or simply marker is a pen which has its own ink source and a tip made of porous, pressed fibers such as felt.
A marker pen consists of a container (glass, aluminum or plastic) and a core of an absorbent material that holds the ink. The upper part of the marker contains the nib that was made in earlier times of a hard felt material, and a cap to prevent the marker from drying out.
Until the early 1990s, the most common solvents that were used for the ink in permanent markers were toluene and xylene. These two substances are both harmful and characterized by a very strong smell. Today, the ink is usually made on the basis of alcohols (e.g. 1-Propanol, 1-butanol, diacetone alcohol and cresols).
Markers may be waterproof, dry-erase, wet-erase (e.g. transparency markers), or permanent.
History
Lee Newman patented a felt-tipped marking pen in 1910. In 1926, Benjamin Paskach patented a "fountain paintbrush", as he called it, which consisted of a sponge-tipped handle containing various paint colors. Markers of this sort began to be popularized with the sale of Sidney Rosenthal's Magic Marker (1953), which consisted of a glass tube of ink with a felt wick. By 1958, use of felt-tipped markers was commonplace for a variety of applications such as lettering, labeling, and creating posters. The year 1962 brought the development of the modern fiber-tipped pen (in contrast to the marker, which generally has a thicker point) by Yukio Horie of the Tokyo Stationery Company (which later became Pentel).
In 1987 Copic Sketch markers were released, further popularising markers for professional illustration.
Parts
The marker reservoir, which holds the ink, is formed from polyester. The "felt" used for the tip is usually made of highly compressed synthetic fibers or porous ceramics. Toluol and xylol were used as solvents for the dye and are still used for the indelible ink in permanent markers. Due to their toxicity, they have often been replaced with less critical substances such as alkyl or cyclic alkylene carbonates (like propylene carbonate) in other types of markers. Water content of the ink can be up to 10%. Besides solvents and the dye itself, the ink may contain additives (e.g. nonylphenylpolyglycol ether, alkylpoly-glycol ether, fatty acid polyglycol ester, or fatty alcohol ethoxalates) and preservatives (e.g. 2-Phenylphenol and its sodium salt, 6-acetoxy-2,4-dimethyl-m-dioxane).
Types
Permanent marker
Permanent markers are porous pens that can write on surfaces such as glass, plastic, wood, metal, and stone. On most surfaces, the ink is generally resistant to rubbing and water, and can last for many years. However, on certain plastics like Teflon, polypropylene etc., the marks made by such pens are not permanent and can be erased easily. Depending on the surface and the marker used, the marks can often be removed with either vigorous scrubbing or chemicals such as acetone.
Highlighters
Highlighters are a form of marker used to highlight and cover over existing writing while still leaving the writing readable. They are generally produced in neon colours to allow for colour coding, as well as attract buyers to them.
Whiteboard markers
A whiteboard marker— or dry-erase marker— uses an erasable ink, made to be used on a slick (or matte-finished), non-porous writing surface, for temporary writing with overhead projectors, whiteboards, glass, and the like. They are designed so that the user can easily erase the marks using a damp cloth, tissue, handkerchief, baby wipe, or other easily cleaned or disposable items. Generally, people use fabrics to do so, but others use items like paper, clothing items, some even use their bare hands to wipe it clear. The erasable ink does not contain the toxic chemical compounds xylene and/or toluene as have been used in permanent markers, being less of a risk to being used as a recreational drug.
Wet-wipe markers— or wet-erase markers— are another version that are used on overhead projectors, signboards, whiteboards, and other non-porous surfaces.
Security marker
Special "security" markers, with fluorescent but otherwise invisible inks, are used for marking valuables in case of burglary. The owner of a stolen, but recovered item can be determined by using ultraviolet light to make the writing visible.
Election marker
Marker pens with election ink (an indelible dye and often a photosensitive agent such as silver nitrate) used to mark the finger, and especially the cuticle, of voters in elections in order to prevent electoral fraud such as double voting. The stain stays visible for a week or two and may also be used to assist in vaccinations.
Porous point pen
A porous point pen contains a point that is made of some porous material such as felt or ceramic. Draftsman's pens usually have a ceramic tip since this wears well and does not broaden when pressure is applied while writing.
Dialectal variations
The use of the terms "marker" and "felt-tipped pen" varies significantly among different parts of the world. This is because most English dialects contain words for particular types of marker, often generic brand names, but there are no such terms in widespread international use.
Asia
In some parts of India, water-based felt-tip pens are referred to as "sketch pens" because they are mainly used for sketching and writing on paper or cardboard. The permanent ink felt-tip markers are referred to as just "markers". In Malaysia and Singapore, marker pens are simply called markers. In the Philippines, a marker is commonly referred to as a "Pentel pen", regardless of brand. In Indonesia, a marker pen is referred to as "Spidol". In South Korea and Japan, marker pens are referred to as "sign pens", "name pens", or "felt pens". Also, permanent pens are also referred to as "Magic" (from a famous pen brand name). In Iran, felt-tip pens are referred to as "Magic" or "Highlight" regardless of its brand.
Australia
In Australia, the term "marker" usually refers only to large-tip markers, and the terms "felt-tip" and "felt pen" usually refer only to fine-tip markers. Markers in Australia are often generically called "texta", after a brand name of a type of permanent marker. Some variation in naming convention occurs between the states, for example in Queensland the brand name "nikko" has been commonly adopted.
New Zealand
The generic terms for fine-tipped markers are usually "felt pen" ,"felt tip pen" or "felts". Large permanent markers are called 'vivids' after a popular brand sold there, the Bic Stephens Vivid
South Africa
In South Africa, the term "Koki" is used for both felt pens and markers, by South Africans, as well as the standard "marker".
Canada and United States
In the United States, the word "marker" is used as well as "magic marker", the latter being a genericized trademark. The word "sharpie" is also now used as a genericized trademark; Sharpie is a popular brand of permanent markers used for labelling. Markers are also sometimes referred to as felt-pens or felts in some parts of Canada.
| Technology | Writing tools | null |
478538 | https://en.wikipedia.org/wiki/Particle%20detector | Particle detector | In experimental and applied particle physics, nuclear physics, and nuclear engineering, a particle detector, also known as a radiation detector, is a device used to detect, track, and/or identify ionizing particles, such as those produced by nuclear decay, cosmic radiation, or reactions in a particle accelerator. Detectors can measure the particle energy and other attributes such as momentum, spin, charge, particle type, in addition to merely registering the presence of the particle.
Examples and types
Many of the detectors invented and used so far are ionization detectors (of which gaseous ionization detectors and semiconductor detectors are most typical) and scintillation detectors; but other, completely different principles have also been applied, like Čerenkov light and transition radiation.
Historical examples
Bubble chamber
Wilson cloud chamber (diffusion chamber)
Photographic plate (Nuclear emulsion)
Detectors for radiation protection
The following types of particle detector are widely used for radiation protection, and are commercially produced in large quantities for general use within the nuclear, medical, and environmental fields.
Dosimeter
Electroscope (when used as a portable dosimeter)
Gaseous ionization detector
Geiger counter
Ionization chamber
Proportional counter
Scintillation counter
Semiconductor detector
Commonly used detectors for particle and nuclear physics
Gaseous ionization detector
Ionization chamber
Proportional counter
Multiwire proportional chamber
Drift chamber
Time projection chamber
Micropattern gaseous detector
Geiger–Müller tube
Spark chamber
Solid-state detectors:
Semiconductor detector and variants including CCDs
Silicon Vertex Detector
Solid-state nuclear track detector
Cherenkov detector
Ring-imaging Cherenkov detector (RICH)
Scintillation counter and associated photomultiplier, photodiode, or avalanche photodiode
Lucas cell
Time-of-flight detector
Transition radiation detector
Calorimeter
Microchannel plate detector
Neutron detector
Modern detectors
Modern detectors in particle physics combine several of the above elements in layers much like an onion.
Research particle detectors
Detectors designed for modern accelerators are huge, both in size and in cost. The term counter is often used instead of detector when the detector counts the particles but does not resolve its energy or ionization. Particle detectors can also usually track ionizing radiation (high energy photons or even visible light). If their main purpose is radiation measurement, they are called radiation detectors, but as photons are also (massless) particles, the term particle detector is still correct.
At colliders
At CERN
for the LHC
CMS
ATLAS
ALICE
LHCb
for the LEP
Aleph
Delphi
L3
Opal
for the SPS
The COMPASS Experiment
Gargamelle
NA61/SHINE
At Fermilab
for the Tevatron
CDF
D0
Mu2e
At DESY
for HERA
H1
HERA-B
HERMES
ZEUS
At BNL
for the RHIC
PHENIX
Phobos
STAR
At SLAC
for the PeP-II
BaBar
for the SLC
SLD
At Cornell
for CESR
CLEO
CUSB
At BINP
for the VEPP-2M and VEPP-2000
ND
SND
CMD
for the VEPP-4
KEDR
Others
MECO from UC Irvine
Under construction
For International Linear Collider (ILC)
CALICE (Calorimeter for Linear Collider Experiment)
Without colliders
Antarctic Muon And Neutrino Detector Array (AMANDA)
Cryogenic Dark Matter Search (CDMS)
Super-Kamiokande
XENON
On spacecraft
Alpha Magnetic Spectrometer (AMS)
DAMPE (DArk Matter Particle Explorer)
Fermi Gamma-ray Space Telescope
JEDI (Jupiter Energetic-particle Detector Instrument)
Theoretical Models of Particle Detectors
Beyond their experimental implementations, theoretical models of particle detectors are also of great importance to theoretical physics. These models consider localized non-relativistic quantum systems coupled to a quantum field. They receive the name of particle detectors because when the non-relativistic quantum system is measured in an excited state, one can claim to have detected a particle. The first instance of particle detector models in the literature dates from the 80's, where a particle in a box was introduced by W. G. Unruh in order to probe a quantum field around a black hole. Shortly after, Bryce DeWitt proposed a simplification of the model, giving rise to the Unruh-DeWitt detector model.
Beyond their applications to theoretical physics, particle detector models are related to experimental fields such as quantum optics, where atoms can be used as detectors for the quantum electromagnetic field via the light-matter interaction. From a conceptual side, particle detectors also allow one to formally define the concept of particles without relying on asymptotic states, or representations of a quantum field theory. As M. Scully puts it, from an operational viewpoint one can state that "a particle is what a particle detector detects", which in essence defines a particle as the detection of excitations of a quantum field.
| Physical sciences | Particle physics: General | null |
478962 | https://en.wikipedia.org/wiki/Overfishing | Overfishing | Overfishing is the removal of a species of fish (i.e. fishing) from a body of water at a rate greater than that the species can replenish its population naturally (i.e. the overexploitation of the fishery's existing fish stock), resulting in the species becoming increasingly underpopulated in that area. Overfishing can occur in water bodies of any sizes, such as ponds, wetlands, rivers, lakes or oceans, and can result in resource depletion, reduced biological growth rates and low biomass levels. Sustained overfishing can lead to critical depensation, where the fish population is no longer able to sustain itself. Some forms of overfishing, such as the overfishing of sharks, has led to the upset of entire marine ecosystems. Types of overfishing include growth overfishing, recruitment overfishing, and ecosystem overfishing. Overfishing not only causes negative impacts on biodiversity and ecosystem functioning, but also reduces fish production, which subsequently leads to negative social and economic consequences.
The ability of a fishery to recover from overfishing depends on whether its overall carrying capacity and the variety of ecological conditions are suitable for the recovery. Dramatic changes in species composition can result in an ecosystem shift, where other equilibrium energy flows involve species compositions different from those that had been present before the depletion of the original fish stock. For example, once trout have been overfished, carp might exploit the change in competitive equilibria and take over in a way that makes it impossible for the trout to re-establish a breeding population.
Since the growth of global fishing enterprises after the 1950s, intensive fishing has spread from a few concentrated areas to encompass nearly all fisheries. The scraping of the ocean floor in bottom dragging is devastating to coral, sponges and other slower-growing benthic species that do not recover quickly, and that provide a habitat for commercial fisheries species. This destruction alters the functioning of the ecosystem and can permanently alter species' composition and biodiversity. Bycatch, the collateral capture of unintended species in the course of fishing, is typically returned to the ocean only to die from injuries or exposure. Bycatch represents about a quarter of all marine catch. In the case of shrimp capture, the mass of bycatch is five times larger than that of the shrimp caught.
A report by FAO in 2020 stated that "in 2017, 34 percent of the fish stocks of the world's marine fisheries were classified as overfished". Mitigation options include: Government regulation, removal of subsidies, minimizing fishing impact, aquaculture and consumer awareness.
Scale
Overfishing has stripped many fisheries around the world of their stocks. The United Nations Food and Agriculture Organization estimated in a 2018 report that 33.1% of world fish stocks are subject to overfishing. Significant overfishing has been observed in pre-industrial times. In particular, the overfishing of the western Atlantic Ocean from the earliest days of European colonisation of the Americas has been well documented.
The fraction of fish stocks that are within biologically sustainable levels has exhibited a decreasing trend, from 90% in 1974 to 66.9% in 2015. In contrast, the percentage of stocks fished at biologically unsustainable levels increased from 10% in 1974 to 33.1% in 2015, with the largest increases in the late-1970s and 1980s.
In 2015, maximally sustainably fished stocks (formerly termed fully fished stocks) accounted for 59.9% and underfished stocks for 7% of the total assessed stocks. While the proportion of underfished stocks decreased continuously from 1974 to 2015, the maximally sustainably fished stocks decreased from 1974 to 1989, and then increased to 59.9% in 2015.
In 2015, among the 16 major statistical areas, the Mediterranean and Black Sea had the highest percentage (62.2%) of unsustainable stocks, closely followed by the Southeast Pacific 61.5% and Southwest Atlantic 58.8%. In contrast, the Eastern Central Pacific, Northeast Pacific (Area 67), Northwest Pacific (Area 61), Western Central Pacific and Southwest Pacific had the lowest proportion (13 to 17%) of fish stocks at biologically unsustainable levels.
Daniel Pauly, a fisheries scientist known for pioneering work on the human impacts on global fisheries, has commented:
According to the Secretary General of the 2002 World Summit on Sustainable Development, "Overfishing cannot continue, the depletion of fisheries poses a major threat to the food supply of millions of people."
The fishing down the food web is something that occurs when overfishing arises. Once all larger fish are caught, the fisherman will start to fish the smaller individuals, which would lead to more fish needing to be caught to keep up with demand. This decreases fish populations, as well as genetic diversity of the species, making them more susceptible to disease, and less likely to adapt to their stressors and the environment. Additionally, catching smaller fish leads to breeding of smaller offspring, which can be problematic for fish. In many species, the smaller the female, the less fecund it is, impacting the fish population.
Types
There are three recognized types of biological overfishing: growth overfishing, recruit overfishing, and ecosystem overfishing.
Growth overfishing
Growth overfishing occurs when fish are harvested at an average size that is smaller than the size that would produce the maximum yield per recruit. A recruit is an individual that makes it to maturity, or into the limits specified by a fishery, which are usually size or age. This makes the total yield less than it would be if the fish were allowed to grow to an appropriate size. It can be countered by reducing fishing mortality to lower levels and increasing the average size of harvested fish to a size that will allow maximum yield per recruit.
Recruitment overfishing
Recruitment overfishing happens when the mature adult population (spawning biomass) is depleted to a level where it no longer has the reproductive capacity to replenish itselfthere are not enough adults to produce offspring. Increasing the spawning stock biomass to a target level is the approach taken by managers to restore an overfished population to sustainable levels. This is generally accomplished by placing moratoriums, quotas, and minimum size limits on a fish population.
Ecosystem overfishing
Ecosystem overfishing occurs when the balance of the ecosystem is altered by overfishing. With declines in the abundance of large predatory species, the abundance of small forage type increases causing a shift in the balance of the ecosystem towards smaller fish species.
Examples and evidence for overfishing
Examples of overfishing exist in areas such as the North Sea, the Grand Banks of Newfoundland and the East China Sea. In these locations, overfishing has not only proved disastrous to fish stocks, but also to the fishing communities relying on the harvest. Like other extractive industries such as forestry and hunting, fisheries are susceptible to economic interaction between ownership or stewardship and sustainability, otherwise known as the tragedy of the commons.
Tuna has been caught by the locals in the upper Adriatic for centuries. Increasing fishing prevented the large schools of little tunny from migrating into the Gulf of Trieste. The last major tuna catch was made in 1954 by the fishermen of Santa Croce, Contovello and Barcola.
The Peruvian coastal anchovy fisheries crashed in the 1970s after overfishing and an El Niño season largely depleted the Peruvian anchovetas from its waters. Anchovies were a major natural resource in Peru; indeed, 1971 alone yielded 10.2 million metric tons of anchovies. However, the following five years saw the Peruvian fleet's catch amount to only about four million tons. This was a major loss to Peru's economy.
The collapse of the Atlantic northwest cod fishery off Newfoundland, and the 1992 decision by Canada to impose an indefinite moratorium on the Grand Banks, is a dramatic example of the consequences of overfishing.
The sole fisheries in the Irish Sea, the west English Channel, and other locations have become overfished to the point of virtual collapse, according to the UK government's official Biodiversity Action Plan. The United Kingdom has created elements in this plan to attempt to restore the fishery, but the expanding global human population and the expanding demand for fish has reached a point where demand for food threatens the stability of these fisheries, if not the species' survival.
Many deep sea fish are at risk, such as orange roughy, rattails, sharks and sablefish. Deep sea fish usually grow slowly because of limited food, have slow metabolisms, low reproductive rates, and many do not reach breeding maturity for 30 to 40 years. A fillet of orange roughy at the store is probably at least 50 years old. Most deep sea fish are in international waters, where there are no legal protections. Most of these fish are caught by deep trawlers near seamounts, where they congregate for food. Flash freezing allows the trawlers to work for days at a time, and modern fishfinders target the fish with ease.
Blue walleye became extinct in the Great Lakes in the 1980s. Until the middle of the 20th century, the walleye was a commercially valuable fish, with about a half million tonnes being landed in the period from about 1880 to the late 1950s, when the populations collapsed, apparently through a combination of overfishing, anthropogenic eutrophication, and competition with introduced rainbow smelt.
The World Wide Fund for Nature and the Zoological Society of London jointly issued their "Living Blue Planet Report" on 16 September 2015 which states that there was a dramatic fall of 74% in worldwide stocks of the important scombridae fish such as mackerel, tuna and bonitos between 1970 and 2010, and the global overall "population sizes of mammals, birds, reptiles, amphibians and fish fell by half on average in just 40 years."
Limited supply due to past overfishing of the Pacific bluefin tuna has contributed to occasional astronomical prices. In January 2019, a 278 kilogram (612 pound) tuna sold for 333.6 million yen, or over US$3 million, US$4,900 per pound.
Sharks and rays: The global abundance of oceanic sharks and rays has declined by 71% since 1970, owing to an 18-fold increase in relative fishing pressure. As a consequence, three-quarters of the species comprising this group are now threatened with extinction. A stark example, caught almost entirely on video, was an incident in Hurghada, Egypt on 8 June 2023, in which Russian Vladimir Popov was killed by a tiger shark in an attack which has been attributed to overfishing of the Red Sea.
A study in 2003 found that, as compared with 1950 levels, only a remnant (in some instances, as little as 10%) of all large ocean-fish stocks are left in the seas. These large ocean fish are the species at the top of the food chains (e.g., tuna, cod, among others). This article was subsequently criticized as being fundamentally flawed, although much debate still exists and the majority of fisheries scientists now consider the results irrelevant with respect to large pelagics (the open seas).
In the United States approximately 27% of exploited fish stocks are considered overfished.
In Tasmania, over 50% of major fisheries species, such as the eastern gemfish, the southern rock lobster, southern bulkefin tuna, jack mackerel, or trumpeter, have declined over the past 75 years due to overfishing.
Consequences
Overfishing not only causes negative impacts on biodiversity and ecosystem functioning, but also reduces fish production, which subsequently leads to negative social and economic consequences. According to a 2008 UN report, the world's fishing fleets are losing US$50 billion each year due to depleted stocks and poor fisheries management. The report, produced jointly by the World Bank and the UN Food and Agriculture Organization (FAO), asserts that half the world's fishing fleet could be scrapped with no change in catch. In addition, the biomass of global fish stocks have been allowed to run down to the point where it is no longer possible to catch the amount of fish that could be caught.
Increased incidence of schistosomiasis in Africa has been linked to declines of fish species that eat the snails carrying the disease-causing parasites.
Massive growth of jellyfish populations threaten fish stocks, as they compete with fish for food, eat fish eggs, and poison or swarm fish, and can survive in oxygen depleted environments where fish cannot; they wreak massive havoc on commercial fisheries. Overfishing eliminates a major jellyfish competitor and predator, exacerbating the jellyfish population explosion. Both climate change and a restructuring of the ecosystem have been found to be major roles in an increase in jellyfish population in the Irish Sea in the 1990s.
According to the 2019 Global Assessment Report on Biodiversity and Ecosystem Services published by the Intergovernmental Science-Policy Platform on Biodiversity and Ecosystem Services, overfishing is a primary driver of mass extinction in the world's oceans. A 2021 study published in the journal Nature asserted that the "primary cause" of ocean defaunation is overfishing. Other studies have shown that overfishing has reduced fish and marine mammal biomass by 60% since the 1800s, and is currently driving over one-third of sharks and rays to extinction.
Acceptable levels
The notion of overfishing hinges on what is meant by an "acceptable level" of fishing. More precise biological and bioeconomic terms define acceptable level as follows:
Biological overfishing occurs when fishing mortality has reached a level where the stock biomass has negative marginal growth (reduced rate of biomass growth), as indicated by the red area in the figure. (Fish are being taken out of the water so quickly that the replenishment of stock by breeding slows down. If the replenishment continues to diminish for long enough, replenishment will go into reverse and the population will decrease.)
Economic or bioeconomic overfishing additionally considers the cost of fishing when determining acceptable catches. Under this framework, a fishery is considered to be overfished when catches exceed maximum economic yield where resource rent is at its maximum. Fish are being removed from the fishery so quickly that the profitability of the fishery is sub-optimal. A more dynamic definition of economic overfishing also considers the present value of the fishery using a relevant discount rate to maximise the flow of resource rent over all future catches.
Harvest control rule
A model proposed in 2010 for predicting acceptable levels of fishing is the Harvest Control Rule (HCR), which is a set of tools and protocols with which management has some direct control of harvest rates and strategies in relation to predicting stock status, and long-term maximum sustainable yields. Constant catch and constant fishing mortality are two types of simple harvest control rules.
Input and output orientations
Fishing capacity can also be defined using an input or output orientation.
An input-oriented fishing capacity is defined as the maximum available capital stock in a fishery that is fully utilized at the maximum technical efficiency in a given time period, given resource and market conditions.
An output-oriented fishing capacity is defined as the maximum catch a vessel (fleet) can produce if inputs are fully utilized given the biomass, the fixed inputs, the age structure of the fish stock, and the present stage of technology.
Technical efficiency of each vessel of the fleet is assumed necessary to attain this maximum catch. The degree of capacity utilization results from the comparison of the actual level of output (input) and the capacity output (input) of a vessel or a fleet.
Reducing overfishing
In order to meet the problems of overfishing, a precautionary approach and Harvest Control Rule (HCR) management principles have been introduced in the main fisheries around the world. The Traffic Light color convention introduces sets of rules based on predefined critical values, which can be adjusted as more information is gained.
The United Nations Convention on the Law of the Sea (UNCLOS) treaty deals with aspects of overfishing in articles 61, 62, and 65.
Article 61 requires all coastal states to ensure that the maintenance of living resources in their exclusive economic zones is not endangered by over-exploitation. The same article addresses the maintenance or restoration of populations of species above levels at which their reproduction may become seriously threatened.
Article 62 provides that coastal states: "shall promote the objective of optimum utilization of the living resources in the exclusive economic zone without prejudice to Article 61"
Article 65 provides generally for the rights of, inter alia, coastal states to prohibit, limit, or regulate the exploitation of marine mammals.
According to some observers, overfishing can be viewed as an example of the tragedy of the commons; appropriate solutions would therefore promote property rights through, for instance, privatization and fish farming. Daniel K. Benjamin, in Fisheries are Classic Example of the 'Tragedy of the Commons', cites research by Grafton, Squires and Fox to support the idea that privatization can solve the overfishing problem: According to recent research on the British Columbia halibut fishery, where the commons has been at least partly privatized, substantial ecological and economic benefits have resulted. There is less damage to fish stocks, the fishing is safer, and fewer resources are needed to achieve a given harvest."
Another possible solution, at least for some areas, is quotas, restricting fishers to a specific quantity of fish. A more radical possibility is declaring certain areas of the sea "no-go zones" and make fishing there strictly illegal, so the fish have time to recover and repopulate.
In order to maximise resources some countries, e.g., Bangladesh and Thailand, have improved the availability of family planning services. The resulting smaller populations have a decreased environmental footprint and reduced food needs.
Controlling consumer behavior and demand is critical in mitigating action. Worldwide, a number of initiatives emerged to provide consumers with information regarding the conservation status of the seafood available to them. The "Guide to Good Fish Guides" lists a number of these.
Government regulation
Many regulatory measures are available for controlling overfishing. These measures include fishing quotas, bag limits, licensing, closed seasons, size limits and the creation of marine reserves and other marine protected areas.
A model of the interaction between fish and fishers showed that when an area is closed to fishers, but there are no catch regulations such as individual transferable quotas, fish catches are temporarily increased but overall fish biomass is reduced, resulting in the opposite outcome from the one desired for fisheries. Thus, a displacement of the fleet from one locality to another will generally have little effect if the same quota is taken. As a result, management measures such as temporary closures or establishing a marine protected area of fishing areas are ineffective when not combined with individual fishing quotas. An inherent problem with quotas is that fish populations vary from year to year. A study has found that fish populations rise dramatically after stormy years due to more nutrients reaching the surface and therefore greater primary production. To fish sustainably, quotas need to be changed each year to account for fish population.
Individual transferable quotas (ITQs) are fishery rationalization instruments defined under the Magnuson-Stevens Fishery Conservation and Management Act as limited access permits to harvest quantities of fish. Fisheries scientists decide the optimal amount of fish (total allowable catch) to be harvested in a certain fishery. The decision considers carrying capacity, regeneration rates and future values. Under ITQs, members of a fishery are granted rights to a percentage of the total allowable catch that can be harvested each year. These quotas can be fished, bought, sold, or leased allowing for the least-cost vessels to be used. ITQs are used in New Zealand, Australia, Iceland, Canada, and the United States.
In 2008, a large-scale study of fisheries that used ITQs compared to ones that did not provide strong evidence that ITQs can help to prevent collapses and restore fisheries that appear to be in decline.
China bans fishing in the South China Sea for a period each year.
Several countries are now effectively managing their fisheries. Examples include Iceland and New Zealand. The United States has turned many of its fisheries around from being in a highly depleted state.
Removal of subsidies
Because government provided financial subsidies can make it economically viable to fish beyond biologically sustainable levels, several scientists have called for an end to fishery subsidies paid to deep-sea fisheries.
Fisheries scientist Daniel Pauly and economist Ussif Rashid Sumaila have examined subsidies paid to bottom trawl fleets around the world. They found that US$152 million per year are paid to deep-sea fisheries. Without these subsidies, global deep-sea fisheries would operate at a loss of US$50 million a year. A great deal of the subsidies paid to deep-sea trawlers is to subsidize the large amount of fuel required to travel beyond the 200 mile limit and drag weighted nets.
"There is surely a better way for governments to spend money than by paying subsidies to a fleet that burns 1.1 billion litres of fuel annually to maintain paltry catches of old growth fish from highly vulnerable stocks, while destroying their habitat in the process" – Pauly.
"Eliminating global subsidies would render these fleets economically unviable and would relieve tremendous pressure on over-fishing and vulnerable deep-sea ecosystems" – Sumaila.Over 30 billion euros in public subsidies are directed to fisheries annually.
Minimizing fishing impact
Fishing techniques may be altered to minimize bycatch and reduce impacts on marine habitats. These techniques include using varied gear types depending on target species and habitat type. For example, a net with larger holes will allow undersized fish to avoid capture. A turtle excluder device (TED) allows sea turtles and other megafauna to escape from shrimp trawls. Avoiding fishing in spawning grounds may allow fish stocks to rebuild by giving adults a chance to reproduce.
Aquaculture
Aquaculture involves the farming of fish in captivity. This approach effectively privatizes fish stocks and creates incentives for farmers to conserve their stocks. It also reduces environmental impact. However, farming carnivorous fish, such as salmon, does not always reduce pressure on wild fisheries, since carnivorous farmed fish are usually fed fishmeal and fish oil extracted from wild forage fish. The various species of Pacific salmon and Atlantic salmon are relatively easy to raise in captivity and such aquacultural operations have existed for more than 150 years. Large scale releases of salmon raised in captivity to supplement wild salmon runs will usually increase fishing pressure on the much less abundant wild salmon runs.
Aquaculture played a minor role in the harvesting of marine organisms until the 1970s. Growth in aquaculture increased rapidly in the 1990s when the rate of wild capture plateaued. Aquaculture now provides approximately half of all harvested aquatic organisms. Aquaculture production rates continue to grow while wild harvest remains steady.
Fish farming can enclose the entire breeding cycle of the fish, with fish being bred in captivity. Some fish prove difficult to breed in captivity and can be caught in the wild as juveniles and brought into captivity to increase their weight. With scientific progress, more species are being made to breed in captivity. This was the case with southern bluefin tuna, which were first bred in captivity in 2009.
Consumer awareness
As global citizens become more aware of overfishing and the ecological destruction of the oceans, movements have sprung up to encourage abstinence—not eating any seafood—or eating only "sustainable seafood".
Sustainable seafood is a movement that has gained momentum as more people become aware of overfishing and environmentally destructive fishing methods. Sustainable seafood is seafood from either fished or farmed sources that can maintain or increase production in the future without jeopardizing the ecosystems from which it was acquired. In general, slow-growing fish that reproduce late in life, such as orange roughy, are vulnerable to overfishing. Seafood species that grow quickly and breed young, such as anchovies and sardines, are much more resistant to overfishing. Several organizations, including the Marine Stewardship Council (MSC), and Friend of the Sea, certify seafood fisheries as sustainable.
The Marine Stewardship Council has developed an environmental standard for sustainable and well-managed fisheries. Environmentally responsible fisheries management and practices are rewarded with the use of its blue product ecolabel. Consumers concerned about overfishing and its consequences are increasingly able to choose seafood products that have been independently assessed against the MSC's environmental standard. This enables consumers to play a part in reversing the decline of fish stocks. As of February 2012, over 100 fisheries around the world have been independently assessed and certified as meeting the MSC standard. Their where-to-buy page lists the currently available certified seafood. As of February 2012, over 13,000 MSC-labelled products are available in 74 countries around the world. is an MSC project to teach schoolchildren about marine environmental issues, including overfishing.
The Monterey Bay Aquarium's Seafood Watch Program, although not an official certifying body like the MSC, also provides guidance on the sustainability of certain fish species. Some seafood restaurants have begun to offer more sustainable seafood options. The Seafood Choices Alliance is an organization whose members include chefs that serve sustainable seafood at their establishments. In the US, the Sustainable Fisheries Act defines sustainable practices through national standards. Although there is no official certifying body like the MSC, the National Oceanic and Atmospheric Administration has created FishWatch to help guide concerned consumers to sustainable seafood choices.
In September 2016, a partnership of Google and Oceana and Skytruth introduced Global Fishing Watch, a website designed to assist citizens of the globe in monitoring fishing activities.
Global goals
The United Nations has included sustainable fishing and ending subsidies that contribute to overfishing as key targets for 2030 as part of their Sustainable Development Goal 14 called "Life Below Water".
Barriers to reducing overfishing
Tragedy of the commons
In open-access resources like fish stocks, in the absence of a system like individual transferable quotas, the impossibility of excluding others provokes the fishermen who want to increase catch to do so effectively by taking someone else' share, intensifying competition. This tragedy of the commons provokes a capitalization process that leads them to increase their costs until they are equal to their revenue, dissipating their rent completely. Causes of the fishing problem can be found in the property rights regime of fishing resources. Overexploitation and rent dissipation of fishermen arise in open-access fisheries as was shown in Gordon.
The fishing industry has a strong financial incentive to oppose some measures aimed at improving the sustainability of fish stocks. Recreational fisherman also have an interest in maintaining access to fish stocks. This leads to extensive lobbying that can block or weaken government policies intended to prevent overfishing.
There is always disagreement between fishermen and government scientists... Imagine an overfished area of the sea in the shape of a hockey field with nets at either end. The few fish left therein would gather around the goals because fish like structured habitats. Scientists would survey the entire field, make lots of unsuccessful hauls, and conclude that it contains few fish. The fishermen would make a beeline to the goals, catch the fish around them, and say the scientists do not know what they are talking about. The subjective impression the fishermen get is always that there's lots of fish - because they only go to places that still have them... fisheries scientists survey and compare entire areas, not only the productive fishing spots. – Fisheries scientist Daniel Pauly
Fish are highly transitory and many species will freely move through different jurisdictions. The conservation efforts of one country can then be exploited by another. Tragedy of the commons can result in perverse incentives to increase fisheries subsidy.
Illegal fishing
While governments can create regulations to control people's behaviors this can be undermined by illegal fishing activity. Estimates of the size of the illegal catch range from 11 to 26 million tonnes, which represents 14-33% of the world's reported catch. Illegal fishing can take many forms. In some developing countries, large numbers of poor people are dependent on fishing. It can prove difficult to regulate this kind of overfishing, especially for weak governments. Even in regulated environments, illegal fishing may occur. While industrial fishing is often effectively controlled, smaller scale and recreational fishermen can often break regulations such as bag limits and seasonal closures. Fisherman can also easily fish illegally by doing things such as underreporting the amount of fish they caught or reporting that they caught one type of fish while actually catching another. There is also a large problem with the surveillance of illegal fishing activity. In 2001, the UN Food and Agriculture Organization (FAO), passed the International Plan of Action to Prevent, Deter and Eliminate Illegal, Unreported and Unregulated Fishing (IPOA-IUU). This is an agreement with the intention to stop port states from allowing boats to dock that participated in illegal, unreported or unregulated fishing. It also gives details for port states on effective measures of inspecting and reporting illegal fishing. Some illegal fishing takes place on an industrial scale.
Territorial disputes
In waters that are the subject of territorial disputes, countries may actively encourage overfishing. A notable example is the cod wars where Britain used its navy to protect its trawlers fishing in Iceland's exclusive economic zone.
International waters
Outside of countries' exclusive economic zones, fishing is difficult to control. Large oceangoing fishing boats are free to exploit fish stocks at will. China is claimed to operate the largest fishing fleet in international waters.
| Physical sciences | Earth science basics: General | Earth science |
479319 | https://en.wikipedia.org/wiki/Japanese%20macaque | Japanese macaque | The Japanese macaque (Macaca fuscata), also known as the snow monkey, is a terrestrial Old World monkey species that is native to Japan. Colloquially, they are referred to as "snow monkeys" because some live in areas where snow covers the ground for months each year – no other non-human primate lives farther north, nor in a colder climate. Individuals have brownish grey fur, pinkish-red faces, and short tails. Two subspecies are known.
In Japan, the species is known as Nihonzaru (ニホンザル, a combination of Nihon 日本 "Japan" + saru 猿 "monkey") to distinguish it from other primates, but the Japanese macaque is very familiar in Japan—as it is the only species of monkey in Japan—so when Japanese people simply say saru, they usually have the Japanese macaque in mind.
Physical characteristics
The Japanese macaque is sexually dimorphic. Males weigh on average , while females average . Macaques from colder areas tend to weigh more than ones from warmer areas. The average height for males is , while the average female height is . The weight of their brain is approximately . Japanese macaques have short stumps for tails that average in males and in females. The macaque has a pinkish face and posterior. The rest of its body is covered in brown or greyish hair. The coat of the macaque is well-adapted to the cold and its thickness increases as temperatures decrease. The macaque can cope with temperatures as low as −20 °C (−4 °F).
Macaques mostly move on all fours. They are semiterrestrial, with females spending more time in the trees and males spending more time on the ground. Macaques are known to leap. They are very good swimmers and have been reported to swim a distance of more than half a kilometer. The lifespan of Japanese macaques is up to 32 years for females and up to 28 years for males, which is high when compared to what typically is seen in other macaque species.
Behavior
Group structure
Japanese macaques live in matrilineal societies, and females stay in their natal groups for life, while males move out before they are sexually mature. Macaque groups tend to contain adults of both sexes. In addition, a Japanese macaque troop contains several matrilines. These matrilines may exist in a dominance hierarchy with all members of a specific group ranking over members of a lower-ranking group. Temporary all-male groups also exist, composed of those who have recently left their natal groups and are about to transfer to another group. However, many males spend ample time away from any group, and may leave and join several groups.
Females of the troop exist in a stable dominance hierarchy and a female's rank depends on that of her mother. Younger females tend to rank higher than their older siblings. Higher-ranking matrilines have greater social cohesion. Strong relationships with dominant females can allow dominant males to retain their rank when they otherwise would not. Males within a group normally have a dominance hierarchy, with one male having alpha status. The dominance status of male macaques usually changes when a former alpha male leaves or dies. Other ways in which status of male hierarchy changes, is when an alpha male loses his rank or when a troop splits, leaving a new alpha male position open. The longer a male is in a troop, the higher his status is likely to be.
Females typically maintain both social relationships and hygiene through grooming. Grooming occurs regardless of climate or season. Females who are matrilineally related groom each other more often than unrelated individuals. Females will groom unrelated females to maintain group cohesion and social relationships between different kinships in a troop. Nevertheless, a female will only groom a limited number of other females, even if the group expands. Females will groom males, usually for hygienic purposes, but that behavior also may serve to attract dominant males to the group. Mothers pass their grooming techniques to their offspring, most probably through social rather than genetic means, as a cultural characteristic.
Documented female troop leadership
Yakei is a female who rose to leadership of her troop at Takasakiyama Natural Zoological Garden in 2021. Her troop consists of 677 Japanese macaque monkeys who live in a sanctuary that was established in 1952 at the zoological garden. At age nine, she overthrew the dominant males in her troop and displaced her high-ranking mother as well. She became the first female leader of the troop during its recorded history of seventy years. Yakei has retained her leadership position through her first breeding season that had been thought to be a time when she might have been challenged successfully. Both scientific and popular interest is leading to extensive coverage of Yakei's behavior.
Mating and parenting
A male and female macaque form a pair bond and mate, feed, rest, and travel together during the mating season, and on average, this relationship typically lasts 16 days. Females enter into consortships with an average of four males a season. Higher-ranking males have longer consortships than their subordinates. In addition, higher-ranking males try to disrupt consortships of lower-ranking males. Females may choose to mate with males of any rank. However, dominant males mate more frequently than others, as they are more successful in mate guarding. The female decides whether mating takes place. In addition, a dominant position does not mean a male will successfully mate with a female. Males may join other troops temporarily during the mating season and mate with those females.
During the mating season, the face and genitalia of males redden and their tails stand erect, and the faces and anogenital regions of females turn scarlet. Macaques copulate both on the ground and in the trees. Roughly one in three copulations leads to ejaculation. Macaques signal when they are ready to mate by looking backward over a shoulder, staying still, or walking backward toward their potential partner. A female emits a "squawk", a "squeak", or produces an atonal "cackle" during copulation. Males have no copulatory vocalizations.
Females engage in same-sex mounting unrelated to the mating season and therefore, are mounted more often by other females than by males. This behavior has led to proposals in literature that female Japanese macaques are generally bisexual, rather than preferentially homo- or heterosexual.
A macaque mother moves to the periphery of her troop to give birth in a secluded spot, unless the group is moving, when the female must stay with it. Macaques usually give birth on the ground. Infants are born with dark-brown hair. A mother and her infant tend to avoid other troop members. The infants consume their first solid food at five to six weeks old, and by seven weeks, can forage independently from their mothers. A mother carries her infant on her belly for its first four weeks. After this time, the mother carries her infant on her back, as well. Infants continue to be carried past a year. The mother may socialize again very slowly. However, alloparenting has been observed, usually by females who have not had infants of their own. Male care of infants occurs in some groups, but not in others; when they do, usually, older males protect, groom, and carry an infant as a female would.
Infants have fully developed their locomotive abilities within three to four months. When an infant is seven months old, its mother discourages suckling; full weaning happens by its eighteenth month.
In some populations, male infants tend to play in larger groups more often than females. However, female infants have more social interaction than their male counterparts, and female infants will associate with individuals of all ages and sexes. When males are two years old, they prefer to associate with other males around the same age.
Communication
During feeding or moving, Japanese macaques often emit sounds that are called "coos". These vocalizations most likely serve to keep the troop together and strengthen social relations among females. Macaques usually respond to coos with coos of their own. Coos also are uttered before grooming along with vocalizations identified as "girney" calls. Variants of the "girney" calls are made in different contexts. This call also serves as appeasement between individuals in aggressive encounters. Macaques have alarm calls for alerting to danger and other calls to signal estrus that sound similar to danger alerts. Threat calls are heard during aggressive encounters and are often uttered by supporters of those involved in antagonistic interactions. The individual being supported supports those callers in the future.
Intelligence and culture
The Japanese macaque is an intelligent species. Researchers studying this species at Koshima Island in Japan left sweet potatoes out on the beach for them to eat, then witnessed one female, named Imo (Japanese for yam or potato), washing the food off with river water rather than brushing it off as the others were doing, and later even dipping her clean food into salty seawater. After a while, other members of her troop started to copy her behavior. This trait was then passed on from generation to generation, until eventually all except the oldest members of the troop were washing their food and even seasoning it in the sea. Similarly, she was the first observed balling up wheat with air pockets and soil, throwing it all into the water, and waiting for the wheat to float back up free from the soil to consume it. An altered misaccount of this incident is the basis for the "hundredth monkey" effect. That behavior also spread among her troop members.
The macaque has other unusual behaviours, including bathing together in hot springs and rolling snowballs for fun. In winter, the bathing is associated with lower levels of stress, with the higher ranking females dominating the restricted resource to compensate for the higher rates of stress outside of the spring. Also, in recent studies, the Japanese macaque has been found to develop different accents, similar to human cultures. Macaques in areas separated by only a few hundred miles may have very different pitches in their calls, their form of communication. The Japanese macaque has been involved in many studies concerning neuroscience and also is used in drug testing.
Ecology
The Japanese macaque is diurnal. In colder areas, from autumn to early winter, macaques feed in between different activities. In the winter, macaques have two to four feeding bouts each day, with fewer daily activities. In the spring and summer, they have two or three bouts of feeding daily. In warmer areas such as Yakushima, daily activities are more varied. The typical day for a macaque is 20.9% inactive, 22.8% traveling, 23.5% feeding, 27.9% social grooming, 1.2% self-grooming, and 3.7% other activities. Macaques usually sleep in trees, but they also sleep on the ground, as well as on or near rocks and fallen trees. During the winter, macaques huddle together for warmth on sleeping grounds. Macaques at Jigokudani Monkey Park are notable for visiting the hot springs in the winter to warm up after being encouraged to concentrate there in the 1960s, part of a plan to reduce local crop damage from foraging.
Diet
The Japanese macaque is omnivorous and eats a variety of foods. More than 213 species of plants are included in the macaque's diet. They also eat insects, bark, and soil. On Yakushima Island, fruit, mature leaves, and fallen seeds are primarily eaten. The macaque also eats fungi, ferns, invertebrates, and other parts of plants. In addition, in Yakushima, their diets vary seasonally with fruits being eaten in the summer and herbs being eaten in the winter. Farther north, macaques mostly eat seasonal foods such as fruit and nuts to store fat for the winter, when food is scarce. On the northern island of Kinkasan, macaques mostly eat fallen seeds, herbs, young leaves, and fruits. When preferred food items are not available, macaques dig up underground plant parts (roots or rhizomes) or eat soil and fish.
Distribution and habitat
The Japanese macaque is the northernmost-living non-human primate. It is found on three of the four main Japanese islands, south of the Blakiston's Line: Honshu, Shikoku, and Kyushu. The northernmost populations live on the Shimokita Peninsula, the northernmost point of Honshu. Several of Japan's smaller islands are inhabited by macaques as well. The southernmost population living on Yakushima Island is a subspecies of the mainland macaques, M. fuscata yakui. A study in 1989 estimated the total population of wild Japanese macaques to be 114,431 individuals.
The Japanese macaque lives in a variety of habitats. It inhabits subtropical forests in the southern part of its range and subarctic forests in mountainous areas in the northern part of its range. It can be found in both warm and cool forests, such as the deciduous forests of central and northern Japan and the broadleaf evergreen forests in the southwest of the islands. Warm temperate evergreen and broadleaf forests and cool temperate deciduous broadleaf forests are the most important habitats for macaques.
In 1972, a troop of approximately 150 Japanese macaques was relocated from Kyoto to a primate observatory in southwest Texas, United States. The observatory is an enclosed ranch-style environment and the macaques have been allowed to roam with minimal human interference. At first, many perished in the unfamiliar habitat, which consists of arid brushland. The macaques eventually adapted to the environment, learned to avoid predators (such as eagles, coyotes, and rattlesnakes), and they learned to forage for mesquite beans, cactus fruits, and other foods. The surviving macaques flourished, and by 1995, the troop consisted of 500 to 600 individuals. In 1996, hunters maimed or killed four escaped macaques; as a result, legal restrictions were publicly clarified and funds were raised to establish a new sanctuary near Dilley, Texas. In 1999, the Animal Protection Institute took over management of the sanctuary and began to rescue other species of primates. As of 2017, the troop cohabitated with six other species of macaque.
Relationship with humans
Traditional human behaviors that are threats to macaques have been slash-and-burn agriculture, use of forest woods for construction and fuel, and hunting. Since World War II, these threats have declined due to social and economic changes in Japan, including the prohibition of macaque hunting in 1947. New threats have emerged, in particular the replacement of natural forests with lumber plantations.
Protection for the macaques, increased afforestation, and the human-caused extinction of their natural predators the Japanese wolf, have led to the macaque population growing heavily since the 1940s. Because of this, and land-use changes increasing the proximity of agriculture to the macaques' range, they have become a major agricultural pest; they can climb over regular fences and quickly realise that deterrents such as scarecrows do not pose an actual threat, so methods such as electric fences must be used. In 2019, the cost of agricultural damage caused by macaques was around 900 million yen. Over 20,000 macaques are culled each year in an attempt to reduce agricultural damage, and there are concerns that this culling is reducing the macaques' range.
Macaques have often entered urban areas, with one macaque recorded living in central Tokyo for several months. In 2022 the city of Yamaguchi experienced aggression from the monkeys with at least 50 people attacked.
Cultural depictions
The Japanese macaque has featured prominently in the religion, folklore, and art of Japan, as well as in proverbs and idiomatic expressions in the Japanese language.
In Shinto belief, mythical beasts known as raijū sometimes appeared as monkeys and kept Raijin, the god of lightning, company. The "three wise monkeys", who warn people to "see no evil, hear no evil and speak no evil", are carved in relief over the door of the famous Tōshō-gū shrine in Nikkō.
The Japanese macaque is a feature of several fairy tales, such as the tale of Momotarō and the fable about The Crab and the Monkey.
The monkey is part of the Chinese zodiac. That zodiac has been used for centuries in Japan and led to many representations of the macaque for that figure.
The creature was sometimes portrayed in paintings of the rich cultural epoch, the Edo period that flourished from 1603 to 1867, as a tangible metaphor for a particular year. The early nineteenth-century artist and samurai, Watanabe Kazan (1793–1841), created a painting of a macaque. The last great master of the ukiyo-e genre of woodblock printing and painting, Tsukioka Yoshitoshi, also featured the macaques in his prints. Also during the Edo period, numerous clasps for kimono or tobacco pouches (collectively called netsuke) were carved in the shape of macaques.
Spoken references to macaques abound in the history of Japan. Originating from before his rise to power, the famed samurai, Toyotomi Hideyoshi, was compared to a monkey in appearance and nicknamed Kozaru ("Little Monkey"). In modern Japanese culture, because monkeys are considered to indulge their libido openly and frequently (much the same way as rabbits are thought to in some Western cultures), a man who is preoccupied with sex might be compared to or metaphorically referred to as a monkey, as might a romantically involved couple who are exceptionally amorous.
| Biology and health sciences | Old World monkeys | Animals |
479373 | https://en.wikipedia.org/wiki/Eraser | Eraser | An eraser (also known as a rubber in some Commonwealth countries, including South Africa from which the material first used got its name) is an article of stationery that is used for removing marks from paper or skin (e.g. parchment or vellum). Erasers have a rubbery consistency and come in a variety of shapes, sizes, and colors. Some pencils have an eraser on one end. Erasers can come in various shapes and colors. Less expensive erasers are made from synthetic rubber and synthetic soy-based gum, but more expensive or specialized erasers are made from vinyl, plastic, or gum-like materials.
At first, erasers were invented to erase mistakes made with a pencil; later, more abrasive ink erasers were introduced. The term is also used for things that remove marks from chalkboards and whiteboards.
History
Before rubber erasers used today, tablets of wax were used to erase lead or charcoal marks from paper. Bits of rough stone such as sandstone or pumice were used to remove small errors from parchment or papyrus documents written in ink. Crustless bread was used; a Meiji period (1868–1912) Tokyo student said: "Bread erasers were used in place of rubber erasers, and so they would give them to us with no restriction on amount. So we thought nothing of taking these and eating a firm part to at least slightly satisfy our hunger."
In 1770 English engineer Edward Nairne is reported to have developed the first widely marketed rubber eraser, for an inventions competition. Until that time the material was known as gum elastic or by its Quechua name (via French) caoutchouc. Nairne sold natural rubber erasers for the high price of three shillings per half-inch cube. According to Nairne, he inadvertently picked up a piece of rubber instead of breadcrumbs, discovered rubber's erasing properties, and began selling rubber erasers. The invention was described by Joseph Priestley on April 15, 1770, in a footnote: "I have seen a substance excellently adapted to the purpose of wiping from paper the mark of black-lead-pencil. ... It is sold by Mr. Nairne, Mathematical Instrument-Maker, opposite the Royal-Exchange." In 1770 the word rubber was in general use for any object used for rubbing; the word became attached to the new material sometime between 1770 and 1778.
However, raw rubber was perishable. In 1839 Charles Goodyear discovered the process of vulcanization, a method that would cure rubber, making it durable. Rubber erasers became common with the advent of vulcanization.
On March 30, 1858, Hymen Lipman of Philadelphia, United States, received the first patent for attaching an eraser to the end of a pencil. It was later invalidated because it was determined to be simply a composite of two devices rather than an entirely new product.
Erasers may be free-standing blocks (block and wedge eraser), or conical caps that can slip onto the end of a pencil (cap eraser). A barrel or click eraser is a device shaped like a pencil, but instead of being filled with pencil lead, its barrel contains a retractable cylinder of eraser material (most commonly soft vinyl). Many, but not all, wooden pencils are made with attached erasers. Novelty erasers made in shapes intended to be amusing are often made of hard vinyl, which tends to smear heavy markings when used as an eraser.
Types
Pencil or cap erasers
Originally made from natural rubber, but now usually from cheaper SBR, this type contains mineral fillers and an abrasive such as pumice with a plasticizer such as vegetable oil. They are relatively hard (in order to remain attached to the pencil) and frequently colored pink.
They can also be permanently attached to the end of a pencil with a ferrule.
Artist's gum eraser
The stylized word "Art gum" was first used in 1903 and trademarked in the United States in 1907. That type of eraser was originally made from oils such as corn oil vulcanized with sulfur dichloride although it may now be made from natural or synthetic rubber or vinyl compounds. It is very soft yet retains its shape and is not mechanically plastic, but crumbles as it is used. It is especially suited to cleaning large areas without damaging the paper. However, they are so soft as to be imprecise in use. The removed graphite is carried away in the crumbles, leaving the eraser clean, but resulting in a lot of eraser residue. This residue must then be brushed away with care, as the eraser particles are coated with the graphite and can make new marks. Art gum erasers are traditionally tan or brown, but some are blue.
Vinyl erasers
High-quality plasticized vinyl or other "plastic" erasers, originally trademarked Mylar in the mid-20th century, are softer, non-abrasive, and erase cleaner than standard rubber erasers. This is because the removed graphite does not remain on the eraser as much as rubber erasers, but is instead absorbed into the discarded vinyl scraps. Being softer and non-abrasive, they are less likely to damage canvas or paper. Engineers favor this type of eraser for work on technical drawings due to their gentleness on paper with less smearing to surrounding areas. They often come in white and can be found in a variety of shapes. More recently, very low-cost erasers are manufactured from highly plasticized vinyl compounds and made in decorative shapes.
Elastomer erasers
In these types, a thermoplastic elastomer combines a styrene resin elastomer and an olefin resin. These erasers have better erasability for erasing pencil marks compared to conventional vinyl erasers. Elastomers can be formed into thin cylindrical or other shapes to be used as extendable erasers.
Kneaded erasers
Kneaded erasers (called putty rubbers outside the United States) have a plastic consistency and are common to most artists' standard toolkit. They can be pulled into a point for erasing small areas and tight detail erasing, molded into a textured surface and used as a reverse stamp to give texture, or used in a "blotting" manner to lighten lines or shading without completely erasing them. They gradually lose their efficacy and resilience as they become infused with particles picked up from erasing and from their environment. They are not suited to erase large areas because of their tendency to deform under vigorous erasing.
Poster putty
Commonly sold in retail outlets with school supplies and home improvement products, this soft, malleable putty appears in many colors and under numerous brand names. Intended to adhere posters and prints to walls without damaging the underlying wall surface, poster putty works much the same as traditional kneaded erasers, but with a greater tack and in some circumstances, lifting strength. Poster putty does not erase so much as lighten by directly pulling particles of graphite, charcoal or pastel from a drawing. In this regard, poster putty does not smudge or damage work in the process. Repeatedly touching the putty to a drawing pulls ever more medium free, gradually lightening the work in a controlled fashion. Poster putty can be shaped into fine points or knife edges, making it ideal for detailed or small areas of work. It can be rolled across a surface to create visual textures. Poster putty loses its efficacy with use, becoming less tacky as the material grows polluted with debris and oils from the user's skin.
Electric erasers
The electric eraser was invented in 1932 by Albert J. Dremel of Racine, Wisconsin, United States. It used a replaceable cylinder of eraser material held by a chuck driven on the axis of a motor. The speed of rotation allowed less pressure to be used, which minimized paper damage. Originally standard pencil-eraser rubber was used, later replaced by higher-performance vinyl. Dremel went on to develop an entire line of hand-held rotary power tools.
Fiberglass erasers
A fiberglass eraser, a bundle of very fine glass fibers, can be used for erasing and other tasks requiring abrasion. Typically the eraser is a pen-shaped device with a replaceable insert with glass fibers, which wear down in use. The fibers are very hard; in addition to removing pencil and pen markings, such erasers are used for cleaning traces on electronic circuit boards to facilitate soldering, removing rust, and many other applications. As an example of an unusual use, a fiberglass eraser was used for preparing a Pterosaur fossil embedded in a very hard and massive limestone. Because fiberglass erasers shed fiberglass dust when used, care must be taken during and after use to avoid accidental contamination with this abrasive dust in sensitive areas of the body, especially in the eyes.
Other
Felt chalkboard erasers or blackboard dusters are used to erase chalk markings on a chalkboard. Chalk writing leaves light-colored particles weakly adhering to a dark surface (e.g., white on black, or yellow on green); it can be rubbed off with a soft material, such as a rag. Erasers for chalkboards are made, with a block of plastic or wood, much larger than an eraser for pen or pencil, with a layer of felt on one side. The block is held in the hand and the felt rubbed against the writing, which it easily wipes off. Chalk dust is released, some of which sticks to the eraser until it is cleaned, usually by hitting it against a hard surface.
Various types of eraser, depending upon the board and the type of ink used, are used to erase a whiteboard.
Dedicated erasers that are supplied with some ballpens and permanent markers are intended only to erase the ink of the writing instrument they are made for; sometimes this is done by making the ink bond more strongly to the material of an eraser than the surface it was applied to.
| Technology | Writing tools | null |
479385 | https://en.wikipedia.org/wiki/Extracellular%20fluid | Extracellular fluid | In cell biology, extracellular fluid (ECF) denotes all body fluid outside the cells of any multicellular organism. Total body water in healthy adults is about 50–60% (range 45 to 75%) of total body weight; women and the obese typically have a lower percentage than lean men. Extracellular fluid makes up about one-third of body fluid, the remaining two-thirds is intracellular fluid within cells. The main component of the extracellular fluid is the interstitial fluid that surrounds cells.
Extracellular fluid is the internal environment of all multicellular animals, and in those animals with a blood circulatory system, a proportion of this fluid is blood plasma. Plasma and interstitial fluid are the two components that make up at least 97% of the ECF. Lymph makes up a small percentage of the interstitial fluid. The remaining small portion of the ECF includes the transcellular fluid (about 2.5%). The ECF can also be seen as having two components – plasma and lymph as a delivery system, and interstitial fluid for water and solute exchange with the cells.
The extracellular fluid, in particular the interstitial fluid, constitutes the body's internal environment that bathes all of the cells in the body. The ECF composition is therefore crucial for their normal functions, and is maintained by a number of homeostatic mechanisms involving negative feedback. Homeostasis regulates, among others, the pH, sodium, potassium, and calcium concentrations in the ECF. The volume of body fluid, blood glucose, oxygen, and carbon dioxide levels are also tightly homeostatically maintained.
The volume of extracellular fluid in a young adult male of 70 kg (154 lbs) is 20% of body weight – about fourteen liters. Eleven liters are interstitial fluid and the remaining three liters are plasma.
Components
The main component of the extracellular fluid (ECF) is the interstitial fluid, or tissue fluid, which surrounds the cells in the body. The other major component of the ECF is the intravascular fluid of the circulatory system called blood plasma. The remaining small percentage of ECF includes the transcellular fluid. These constituents are often called "fluid compartments". The volume of extracellular fluid in a young adult male of 70 kg, is 20% of body weight – about fourteen liters.
Interstitial fluid
Interstitial fluid is essentially comparable to plasma. The interstitial fluid and plasma make up about 97% of the ECF, and a small percentage of this is lymph.
Interstitial fluid is the body fluid between blood vessels and cells, containing nutrients from capillaries by diffusion and holding waste products discharged by cells due to metabolism. 11 liters of the ECF are interstitial fluid and the remaining three liters are plasma. Plasma and interstitial fluid are very similar because water, ions, and small solutes are continuously exchanged between them across the walls of capillaries, through pores and capillary clefts.
Interstitial fluid consists of a water solvent containing sugars, salts, fatty acids, amino acids, coenzymes, hormones, neurotransmitters, white blood cells and cell waste-products. This solution accounts for 26% of the water in the human body. The composition of interstitial fluid depends upon the exchanges between the cells in the biological tissue and the blood. This means that tissue fluid has a different composition in different tissues and in different areas of the body.
The plasma that filters through the blood capillaries into the interstitial fluid does not contain red blood cells or platelets as they are too large to pass through but can contain some white blood cells to help the immune system.
Once the extracellular fluid collects into small vessels (lymph capillaries) it is considered to be lymph, and the vessels that carry it back to the blood are called the lymphatic vessels. The lymphatic system returns protein and excess interstitial fluid to the circulation.
The ionic composition of the interstitial fluid and blood plasma vary due to the Gibbs–Donnan effect. This causes a slight difference in the concentration of cations and anions between the two fluid compartments.
Transcellular fluid
Transcellular fluid is formed from the transport activities of cells, and is the smallest component of extracellular fluid. These fluids are contained within epithelial lined spaces. Examples of this fluid are cerebrospinal fluid, aqueous humor in the eye, serous fluid in the serous membranes lining body cavities, perilymph and endolymph in the inner ear, and joint fluid. Due to the varying locations of transcellular fluid, the composition changes dramatically. Some of the electrolytes present in the transcellular fluid are sodium ions, chloride ions, and bicarbonate ions.
Function
Extracellular fluid provides the medium for the exchange of substances between the ECF and the cells, and this can take place through dissolving, mixing and transporting in the fluid medium. Substances in the ECF include dissolved gases, nutrients, and electrolytes, all needed to maintain life. ECF also contains materials secreted from cells in soluble form, but which quickly coalesce into fibers (e.g. collagen, reticular, and elastic fibres) or precipitates out into a solid or semisolid form (e.g. proteoglycans which form the bulk of cartilage, and the components of bone). These and many other substances occur, especially in association with various proteoglycans, to form the extracellular matrix, or the "filler" substance, between the cells throughout the body. These substances occur in the extracellular space, and are therefore all bathed or soaked in ECF, without being part of it.
Oxygenation
One of the main roles of extracellular fluid is to facilitate the exchange of molecular oxygen from blood to tissue cells and for carbon dioxide, CO2, produced in cell mitochondria, back to the blood. Since carbon dioxide is about 20 times more soluble in water than oxygen, it can relatively easily diffuse in the aqueous fluid between cells and blood.
However, hydrophobic molecular oxygen has very poor water solubility and prefers hydrophobic lipid crystalline structures. As a result of this, plasma lipoproteins can carry significantly more O2 than in the surrounding aqueous medium.
If hemoglobin in erythrocytes is the main transporter of oxygen in the blood, plasma lipoproteins may be its only carrier in the ECF.
The oxygen-carrying capacity of lipoproteins, reduces in ageing and inflammation. This results in changes of ECF functions, reduction of tissue O2 supply and contributes to development of tissue hypoxia. These changes in lipoproteins are caused by oxidative or inflammatory damage.
Regulation
The internal environment is stabilised in the process of homeostasis. Complex homeostatic mechanisms operate to regulate and keep the composition of the ECF stable. Individual cells can also regulate their internal composition by various mechanisms.
There is a significant difference between the concentrations of sodium and potassium ions inside and outside the cell. The concentration of sodium ions is considerably higher in the extracellular fluid than in the intracellular fluid. The converse is true of the potassium ion concentrations inside and outside the cell. These differences cause all cell membranes to be electrically charged, with the positive charge on the outside of the cells and the negative charge on the inside. In a resting neuron (not conducting an impulse) the membrane potential is known as the resting potential, and between the two sides of the membrane is about −70 mV.
This potential is created by sodium–potassium pumps in the cell membrane, which pump sodium ions out of the cell, into the ECF, in return for potassium ions which enter the cell from the ECF. The maintenance of this difference in the concentration of ions between the inside of the cell and the outside, is critical to keep normal cell volumes stable, and also to enable some cells to generate action potentials.
In several cell types voltage-gated ion channels in the cell membrane can be temporarily opened under specific circumstances for a few microseconds at a time. This allows a brief inflow of sodium ions into the cell (driven in by the sodium ion concentration gradient that exists between the outside and inside of the cell). This causes the cell membrane to temporarily depolarize (lose its electrical charge) forming the basis of action potentials.
The sodium ions in the ECF also play an important role in the movement of water from one body compartment to the other. When tears are secreted, or saliva is formed, sodium ions are pumped from the ECF into the ducts in which these fluids are formed and collected. The water content of these solutions results from the fact that water follows the sodium ions (and accompanying anions) osmotically. The same principle applies to the formation of many other body fluids.
Calcium ions have a great propensity to bind to proteins. This changes the distribution of electrical charges on the protein, with the consequence that the 3D (or tertiary) structure of the protein is altered. The normal shape, and therefore function of very many of the extracellular proteins, as well as the extracellular portions of the cell membrane proteins, is dependent on a very precise ionized calcium concentration in the ECF. The proteins that are particularly sensitive to changes in the ECF ionized calcium concentration are several of the clotting factors in the blood plasma, which are functionless in the absence of calcium ions, but become fully functional on the addition of the correct concentration of calcium salts. The voltage gated sodium ion channels in the cell membranes of nerves and muscle have an even greater sensitivity to changes in the ECF ionized calcium concentration. Relatively small decreases in the plasma ionized calcium levels (hypocalcemia) cause these channels to leak sodium into the nerve cells or axons, making them hyper-excitable, thus causing spontaneous muscle spasms (tetany) and paraesthesia (the sensation of "pins and needles") of the extremities and round the mouth. When the plasma ionized calcium rises above normal (hypercalcemia) more calcium is bound to these sodium channels having the opposite effect, causing lethargy, muscle weakness, anorexia, constipation and labile emotions.
The tertiary structure of proteins is also affected by the pH of the bathing solution. In addition, the pH of the ECF affects the proportion of the total amount of calcium in the plasma which occurs in the free, or ionized form, as opposed to the fraction that is bound to protein and phosphate ions. A change in the pH of the ECF therefore alters the ionized calcium concentration of the ECF. Since the pH of the ECF is directly dependent on the partial pressure of carbon dioxide in the ECF, hyperventilation, which lowers the partial pressure of carbon dioxide in the ECF, produces symptoms that are almost indistinguishable from low plasma ionized calcium concentrations.
The extracellular fluid is constantly "stirred" by the circulatory system, which ensures that the watery environment which bathes the body's cells is virtually identical throughout the body. This means that nutrients can be secreted into the ECF in one place (e.g. the gut, liver, or fat cells) and will, within about a minute, be evenly distributed throughout the body. Hormones are similarly rapidly and evenly spread to every cell in the body, regardless of where they are secreted into the blood. Oxygen taken up by the lungs from the alveolar air is also evenly distributed at the correct partial pressure to all the cells of the body. Waste products are also uniformly spread to the whole of the ECF, and are removed from this general circulation at specific points (or organs), once again ensuring that there is generally no localized accumulation of unwanted compounds or excesses of otherwise essential substances (e.g. sodium ions, or any of the other constituents of the ECF). The only significant exception to this general principle is the plasma in the veins, where the concentrations of dissolved substances in individual veins differ, to varying degrees, from those in the rest of the ECF. However, this plasma is confined within the waterproof walls of the venous tubes, and therefore does not affect the interstitial fluid in which the body's cells live. When the blood from all the veins in the body mixes in the heart and lungs, the differing compositions cancel out (e.g. acidic blood from active muscles is neutralized by the alkaline blood homeostatically produced by the kidneys). From the left atrium onward, to every organ in the body, the normal, homeostatically regulated values of all of the ECF's components are therefore restored.
Interaction between the blood plasma, interstitial fluid and lymph
The arterial blood plasma, interstitial fluid and lymph interact at the level of the blood capillaries. The capillaries are permeable and water can move freely in and out. At the arteriolar end of the capillary the blood pressure is greater than the hydrostatic pressure in the tissues. Water will therefore seep out of the capillary into the interstitial fluid. The pores through which this water moves are large enough to allow all the smaller molecules (up to the size of small proteins such as insulin) to move freely through the capillary wall as well. This means that their concentrations across the capillary wall equalize, and therefore have no osmotic effect (because the osmotic pressure caused by these small molecules and ions – called the crystalloid osmotic pressure to distinguish it from the osmotic effect of the larger molecules that cannot move across the capillary membrane – is the same on both sides of capillary wall).
The movement of water out of the capillary at the arteriolar end causes the concentration of the substances that cannot cross the capillary wall to increase as the blood moves to the venular end of the capillary. The most important substances that are confined to the capillary tube are plasma albumin, the plasma globulins and fibrinogen. They, and particularly the plasma albumin, because of its molecular abundance in the plasma, are responsible for the so-called "oncotic" or "colloid" osmotic pressure which draws water back into the capillary, especially at the venular end.
The net effect of all of these processes is that water moves out of and back into the capillary, while the crystalloid substances in the capillary and interstitial fluids equilibrate. Since the capillary fluid is constantly and rapidly renewed by the flow of the blood, its composition dominates the equilibrium concentration that is achieved in the capillary bed. This ensures that the watery environment of the body's cells is always close to their ideal environment (set by the body's homeostats).
A small proportion of the solution that leaks out of the capillaries is not drawn back into the capillary by the colloid osmotic forces. This amounts to between 2–4 liters per day for the body as a whole. This water is collected by the lymphatic system and is ultimately discharged into the left subclavian vein, where it mixes with the venous blood coming from the left arm, on its way to the heart. The lymph flows through lymph capillaries to lymph nodes where bacteria and tissue debris are removed from the lymph, while various types of white blood cells (mainly lymphocytes) are added to the fluid. In addition the lymph which drains the small intestine contains fat droplets called chylomicrons after the ingestion of a fatty meal. This lymph is called chyle which has a milky appearance, and imparts the name lacteals (referring to the milky appearance of their contents) to the lymph vessels of the small intestine.
Extracellular fluid may be mechanically guided in this circulation by the vesicles between other structures. Collectively this forms the interstitium, which may be considered a newly identified biological structure in the body. However, there is some debate over whether the interstitium is an organ.
Electrolytic constituents
Main cations:
Sodium (Na+) 136–146 mM
Potassium (K+) 3.8–5.0 mM
Calcium (Ca2+) 1.0–1.4 mM
Main anions:
Chloride (Cl−) 103–112 mM
Bicarbonate (HCO3−) 22–28 mM
Phosphate (HPO42−) 0.8–1.4 mM
| Biology and health sciences | Animal anatomy and morphology | Biology |
479392 | https://en.wikipedia.org/wiki/Myosin | Myosin | Myosins () are a family of motor proteins (though most often protein complexes) best known for their roles in muscle contraction and in a wide range of other motility processes in eukaryotes. They are ATP-dependent and responsible for actin-based motility.
The first myosin (M2) to be discovered was in 1864 by Wilhelm Kühne. Kühne had extracted a viscous protein from skeletal muscle that he held responsible for keeping the tension state in muscle. He called this protein myosin. The term has been extended to include a group of similar ATPases found in the cells of both striated muscle tissue and smooth muscle tissue.
Following the discovery in 1973 of enzymes with myosin-like function in Acanthamoeba castellanii, a global range of divergent myosin genes have been discovered throughout the realm of eukaryotes.
Although myosin was originally thought to be restricted to muscle cells (hence myo-(s) + -in), there is no single "myosin"; rather it is a very large superfamily of genes whose protein products share the basic properties of actin binding, ATP hydrolysis (ATPase enzyme activity), and force transduction. Virtually all eukaryotic cells contain myosin isoforms. Some isoforms have specialized functions in certain cell types (such as muscle), while other isoforms are ubiquitous. The structure and function of myosin is globally conserved across species, to the extent that rabbit muscle myosin II will bind to actin from an amoeba.
Structure and functions
Domains
Most myosin molecules are composed of a head, neck, and tail domain.
The head domain binds the filamentous actin, and uses ATP hydrolysis to generate force and to "walk" along the filament towards the barbed (+) end (with the exception of myosin VI, which moves towards the pointed (-) end).
the neck domain acts as a linker and as a lever arm for transducing force generated by the catalytic motor domain. The neck domain can also serve as a binding site for myosin light chains which are distinct proteins that form part of a macromolecular complex and generally have regulatory functions.
The tail domain generally mediates interaction with cargo molecules and/or other myosin subunits. In some cases, the tail domain may play a role in regulating motor activity.
Power stroke
Multiple myosin II molecules generate force in skeletal muscle through a power stroke mechanism fuelled by the energy released from ATP hydrolysis. The power stroke occurs at the release of phosphate from the myosin molecule after the ATP hydrolysis while myosin is tightly bound to actin. The effect of this release is a conformational change in the molecule that pulls against the actin. The release of the ADP molecule leads to the so-called rigor state of myosin. The binding of a new ATP molecule will release myosin from actin. ATP hydrolysis within the myosin will cause it to bind to actin again to repeat the cycle. The combined effect of the myriad power strokes causes the muscle to contract.
Nomenclature, evolution, and the family tree
The wide variety of myosin genes found throughout the eukaryotic phyla were named according to different schemes as they were discovered. The nomenclature can therefore be somewhat confusing when attempting to compare the functions of myosin proteins within and between organisms.
Skeletal muscle myosin, the most conspicuous of the myosin superfamily due to its abundance in muscle fibers, was the first to be discovered. This protein makes up part of the sarcomere and forms macromolecular filaments composed of multiple myosin subunits. Similar filament-forming myosin proteins were found in cardiac muscle, smooth muscle, and nonmuscle cells. However, beginning in the 1970s, researchers began to discover new myosin genes in simple eukaryotes encoding proteins that acted as monomers and were therefore entitled Class I myosins. These new myosins were collectively termed "unconventional myosins" and have been found in many tissues other than muscle. These new superfamily members have been grouped according to phylogenetic relationships derived from a comparison of the amino acid sequences of their head domains, with each class being assigned a Roman numeral (see phylogenetic tree). The unconventional myosins also have divergent tail domains, suggesting unique functions. The now diverse array of myosins likely evolved from an ancestral precursor (see picture).
Analysis of the amino acid sequences of different myosins shows great variability among the tail domains, but strong conservation of head domain sequences. Presumably this is so the myosins may interact, via their tails, with a large number of different cargoes, while the goal in each case – to move along actin filaments – remains the same and therefore requires the same machinery in the motor. For example, the human genome contains over 40 different myosin genes.
These differences in shape also determine the speed at which myosins can move along actin filaments. The hydrolysis of ATP and the subsequent release of the phosphate group causes the "power stroke", in which the "lever arm" or "neck" region of the heavy chain is dragged forward. Since the power stroke always moves the lever arm by the same angle, the length of the lever arm determines the displacement of the cargo relative to the actin filament. A longer lever arm will cause the cargo to traverse a greater distance even though the lever arm undergoes the same angular displacement – just as a person with longer legs can move farther with each individual step. The velocity of a myosin motor depends upon the rate at which it passes through a complete kinetic cycle of ATP binding to the release of ADP.
Myosin classes
Myosin I
Myosin I, a ubiquitous cellular protein, functions as monomer and functions in vesicle transport. It has a step size of 10 nm and has been implicated as being responsible for the adaptation response of the stereocilia in the inner ear.
Myosin II
Myosin II (also known as conventional myosin) is the myosin type responsible for producing muscle contraction in muscle cells in most animal cell types. It is also found in non-muscle cells in contractile bundles called stress fibers.
Myosin II contains two heavy chains, each about 2000 amino acids in length, which constitute the head and tail domains. Each of these heavy chains contains the N-terminal head domain, while the C-terminal tails take on a coiled-coil morphology, holding the two heavy chains together (imagine two snakes wrapped around each other, as in a caduceus). Thus, myosin II has two heads. The intermediate neck domain is the region creating the angle between the head and tail. In smooth muscle, a single gene (MYH11)) codes for the heavy chains myosin II, but splice variants of this gene result in four distinct isoforms.
It also contains 4 myosin light chains (MLC), resulting in 2 per head, weighing 20 (MLC20) and 17 (MLC17) kDa. These bind the heavy chains in the "neck" region between the head and tail.
The MLC20 is also known as the regulatory light chain and actively participates in muscle contraction.
The MLC17 is also known as the essential light chain. Its exact function is unclear, but is believed to contribute to the structural stability of the myosin head along with MLC20. Two variants of MLC17 (MLC17a/b) exist as a result of alternative splicing at the MLC17 gene.
In muscle cells, the long coiled-coil tails of the individual myosin molecules can auto-inhibit active function in the 10S conformation or upon phosphorylation, change to the 6S conformation and join, forming the thick filaments of the sarcomere. The force-producing head domains stick out from the side of the thick filament, ready to walk along the adjacent actin-based thin filaments in response to the proper chemical signals and may be in either auto-inhibited or active conformation. The balance/transition between active and inactive states is subject to extensive chemical regulation.
Myosin III
Myosin III is a poorly understood member of the myosin family. It has been studied in vivo in the eyes of Drosophila, where it is thought to play a role in phototransduction. A human homologue gene for myosin III, MYO3A, has been uncovered through the Human Genome Project and is expressed in the retina and cochlea.
Myosin IV
Myosin IV has a single IQ motif and a tail that lacks any coiled-coil forming sequence. It has homology similar to the tail domains of Myosin VII and XV.
Myosin V
Myosin V is an unconventional myosin motor, which is processive as a dimer and has a step size of 36 nm. It translocates (walks) along actin filaments traveling towards the barbed end (+ end) of the filaments. Myosin V is involved in the transport of cargo (e.g. RNA, vesicles, organelles, mitochondria) from the center of the cell to the periphery, but has been furthermore shown to act like a dynamic tether, retaining vesicles and organelles in the actin-rich periphery of cells. A recent single molecule in vitro reconstitution study on assembling actin filaments suggests that Myosin V travels farther on newly assembling (ADP-Pi rich) F-actin, while processive runlengths are shorter on older (ADP-rich) F-actin.
The Myosin V motor head can be subdivided into the following functional regions:
Nucleotide-binding site - These elements together coordinate di-valent metal cations (usually magnesium) and catalyze hydrolysis:
Switch I - This contains a highly conserved SSR motif. Isomerizes in the presence of ATP.
Switch II - This is the Kinase-GTPase version of the Walker B motif DxxG. Isomerizes in the presence of ATP.
P-loop - This contains the Walker A motif GxxxxGK(S,T). This is the primary ATP binding site.
Transducer - The seven β-strands that underpin the motor head's structure.
U50 and L50 - The Upper (U50) and Lower (L50) domains are each around 50kDa. Their spatial separation forms a cleft critical for binding to actin and some regulatory compounds.
SH1 helix and Relay - These elements together provide an essential mechanism for coupling the enzymatic state of the motor domain to the powerstroke-producing region (converter domain, lever arm, and light chains).
Converter - This converts a change of conformation in the motor head to an angular displacement of the lever arm (in most cases reinforced with light chains).
Myosin VI
Myosin VI is an unconventional myosin motor, which is primarily processive as a dimer, but also acts as a nonprocessive monomer. It walks along actin filaments, travelling towards the pointed end (- end) of the filaments. Myosin VI is thought to transport endocytic vesicles into the cell.
Myosin VII
Myosin VII is an unconventional myosin with two FERM domains in the tail region. It has an extended lever arm consisting of five calmodulin binding IQ motifs followed by a single alpha helix (SAH) Myosin VII is required for phagocytosis in Dictyostelium discoideum, spermatogenesis in C. elegans and stereocilia formation in mice and zebrafish.
Myosin VIII
Myosin VIII is a plant-specific myosin linked to cell division; specifically, it is involved in regulating the flow of cytoplasm between cells and in the localization of vesicles to the phragmoplast.
Myosin IX
Myosin IX is a group of single-headed motor proteins. It was first shown to be minus-end directed, but a later study showed that it is plus-end directed. The movement mechanism for this myosin is poorly understood.
Myosin X
Myosin X is an unconventional myosin motor, which is functional as a dimer. The dimerization of myosin X is thought to be antiparallel. This behavior has not been observed in other myosins. In mammalian cells, the motor is found to localize to filopodia. Myosin X walks towards the barbed ends of filaments. Some research suggests it preferentially walks on bundles of actin, rather than single filaments. It is the first myosin motor found to exhibit this behavior.
Myosin XI
Myosin XI directs the movement of organelles such as plastids and mitochondria in plant cells. It is responsible for the light-directed movement of chloroplasts according to light intensity and the formation of stromules interconnecting different plastids. Myosin XI also plays a key role in polar root tip growth and is necessary for proper root hair elongation. A specific Myosin XI found in Nicotiana tabacum was discovered to be the fastest known processive molecular motor, moving at 7μm/s in 35 nm steps along the actin filament.
Myosin XII
Myosin XIII
Myosin XIV
This myosin group has been found in the Apicomplexa phylum. The myosins localize to plasma membranes of the intracellular parasites and may then be involved in the cell invasion process.
This myosin is also found in the ciliated protozoan Tetrahymena thermophila. Known functions include: transporting phagosomes to the nucleus and perturbing the developmentally regulated elimination of the macronucleus during conjugation.
Myosin XV
Myosin XV is necessary for the development of the actin core structure of the non-motile stereocilia located in the inner ear. It is thought to be functional as a monomer.
Myosin XVI
Myosin XVII
Myosin XVIII
MYO18A A gene on chromosome 17q11.2 that encodes actin-based motor molecules with ATPase activity, which may be involved in maintaining stromal cell scaffolding required for maintaining intercellular contact.
Myosin XIX
Unconventional myosin XIX (Myo19) is a mitochondrial associated myosin motor.
Genes in humans
Note that not all of these genes are active.
Class I: MYO1A, MYO1B, MYO1C, MYO1D, MYO1E, MYO1F, MYO1G, MYO1H
Class II: MYH1, MYH2, MYH3, MYH4, MYH6, MYH7, MYH7B, MYH8, MYH9, MYH10, MYH11, MYH13, MYH14, MYH15, MYH16
Class III: MYO3A, MYO3B
Class V: MYO5A, MYO5B, MYO5C
Class VI: MYO6
Class VII: MYO7A, MYO7B
Class IX: MYO9A, MYO9B
Class X: MYO10
Class XV: MYO15A, MYO15B
Class XVI: MYO16
Class XVIII: MYO18A, MYO18B
Class XIX: MYO19
Myosin light chains are distinct and have their own properties. They are not considered "myosins" but are components of the macromolecular complexes that make up the functional myosin enzymes.
Light chain: MYL1, MYL2, MYL3, MYL4, MYL5, MYL6, MYL6B, MYL7, MYL9, MYLIP, MYLK, MYLK2, MYLL1
Paramyosin
Paramyosin is a large, 93-115kDa muscle protein that has been described in a number of diverse invertebrate phyla. Invertebrate thick filaments are thought to be composed of an inner paramyosin core surrounded by myosin. The myosin interacts with actin, resulting in fibre contraction. Paramyosin is found in many different invertebrate species, for example, Brachiopoda, Sipunculidea, Nematoda, Annelida, Mollusca, Arachnida, and Insecta. Paramyosin is responsible for the "catch" mechanism that enables sustained contraction of muscles with very little energy expenditure, such that a clam can remain closed for extended periods.
Paramyosins can be found in seafood. A recent computational study showed that following human intestinal digestion, paramyosins of common octopus, Humboldt squid, Japanese abalone, Japanese scallop, Mediterranean mussel, Pacific oyster, sea cucumber, and Whiteleg shrimp could release short peptides that inhibit the enzymatic activities of angiotensin converting enzyme and dipeptidyl peptidase.
| Biology and health sciences | Cell parts | Biology |
479444 | https://en.wikipedia.org/wiki/Arteriole | Arteriole | An arteriole is a small-diameter blood vessel in the microcirculation that extends and branches out from an artery and leads to capillaries.
Arterioles have muscular walls (usually only one to two layers of smooth muscle cells) and are the primary site of vascular resistance. The greatest change in blood pressure and velocity of blood flow occurs at the transition of arterioles to capillaries. This function is extremely important because it prevents the thin, one-layer capillaries from exploding upon pressure. The arterioles achieve this decrease in pressure, as they are the site with the highest resistance (a large contributor to total peripheral resistance) which translates to a large decrease in the pressure.
Structure
In a healthy vascular system, the endothelium lines all blood-contacting surfaces, including arteries, arterioles, veins, venules, capillaries, and heart chambers. This healthy condition is promoted by the ample production of nitric oxide by the endothelium, which requires a biochemical reaction regulated by a complex balance of polyphenols, various nitric oxide synthase enzymes and L-arginine. In addition, there is direct electrical and chemical communication via gap junctions between the endothelial cells and the vascular smooth muscle.
Physiology
Blood pressure
Blood pressure in the arteries supplying the body is a result of the work needed to pump the cardiac output (the flow of blood pumped by the heart) through the vascular resistance, sometimes termed total peripheral resistance. An increase in the tunica media to luminal diameter ratio has been observed in hypertensive arterioles (arteriolosclerosis) as the vascular wall thickens and/or luminal diameter decreases.
The up and down fluctuation of the arterial blood pressure is due to the pulsatile nature of the cardiac output and determined by the interaction of the stroke volume versus the volume and elasticity of the major arteries.
The decreased velocity of flow in the capillaries increases the blood pressure, due to Bernoulli's principle. This induces gas and nutrients to move from the blood to the cells, due to the lower osmotic pressure outside the capillary. The opposite process occurs when the blood leaves the capillaries and enters the venules, where the blood pressure drops due to an increase in flow rate. Arterioles receive autonomic nervous system innervation and respond to various circulating hormones in order to regulate their diameter. Retinal vessels lack a functional sympathetic innervation.
Autoregulation
Arteriole diameter varies according to autoregulation of organs or tissues to maintain sufficient blood flow despite changes in pressure via metabolic or myogenic factors which include stretch, carbon dioxide, and oxygen among other factors. Generally, norepinephrine and epinephrine (hormones produced by sympathetic nerves and the adrenal gland medulla) are vasoconstrictive acting on alpha 1-adrenergic receptors. However, the arterioles of skeletal muscle, cardiac muscle, and pulmonary circulation vasodilate in response to these hormones when they act on beta-adrenergic receptors. Generally, stretch and high oxygen tension increase tone, and carbon dioxide and low pH promote vasodilation. Pulmonary arterioles are a noteworthy exception as they vasodilate in response to high oxygen. Brain arterioles are particularly sensitive to pH with reduced pH promoting vasodilation. A number of hormones influence arteriole tone such as angiotensin II (vasoconstrictive), endothelin (vasoconstrictive), bradykinin (vasodilation), atrial natriuretic peptide (vasodilation), and prostacyclin (vasodilation).
Clinical significance
Arteriole diameters decrease with age and with exposure to air pollution.
Disease
Any pathology which constricts blood flow, such as stenosis, will increase total peripheral resistance and lead to hypertension.
Arteriosclerosis
Arteriolosclerosis is the term specifically used for the hardening of arteriole walls. This can be due to decreased elastic production from fibrinogen, associated with ageing, or hypertension or pathological conditions such as atherosclerosis.
Arteritis
Arteritis of the arterioles occurs when the arteriole walls become inflamed as a result of either an immune response to infection or an autoimmune response.
Medication
The muscular contraction of arterioles is targeted by drugs that lower blood pressure (antihypertensives), for example the dihydropyridines (nifedipine and nicardipine), which block the calcium conductance in the muscular layer of the arterioles, causing relaxation.
This decreases the resistance to flow into peripheral vascular beds, lowering overall systemic pressure.
Metarterioles
A "metarteriole" is an arteriole which bypasses capillary circulation.
| Biology and health sciences | Circulatory system | Biology |
479942 | https://en.wikipedia.org/wiki/Asclepias | Asclepias | Asclepias is a genus of herbaceous, perennial, flowering plants known as milkweeds, named for their latex, a milky substance containing cardiac glycosides termed cardenolides, exuded where cells are damaged. Most species are toxic to humans and many other species, primarily due to the presence of cardenolides. However, as with many such plants, some species feed upon milkweed leaves or the nectar from their flowers. A noteworthy feeder on milkweeds is the monarch butterfly, which uses and requires certain milkweeds as host plants for their larvae.
The Asclepias genus contains over 200 species distributed broadly across Africa, North America, and South America. It previously belonged to the family Asclepiadaceae, which is now classified as the subfamily Asclepiadoideae of the dogbane family, Apocynaceae.
The genus was formally described by Carl Linnaeus in 1753, who named it after Asclepius, the Greek god of healing.
Flowers
Members of the genus produce some of the most complex flowers in the plant kingdom, comparable to orchids in complexity. Five petals reflex backwards revealing a gynostegium surrounded by a five-membrane corona. The corona is composed of a five-paired hood-and-horn structure with the hood acting as a sheath for the inner horn. Glands holding pollinia are found between the hoods. The size, shape and color of the horns and hoods are often important identifying characteristics for species in the genus Asclepias.
Pollination in this genus is accomplished in an unusual manner. Pollen is grouped into complex structures called pollinia (or "pollen sacs"), rather than being individual grains or tetrads, as is typical for most plants. The feet or mouthparts of flower-visiting insects, such as bees, wasps, and butterflies, slip into one of the five slits in each flower formed by adjacent anthers. The bases of the pollinia then mechanically attach to the insect, so that a pair of pollen sacs can be pulled free when the pollinator flies off, assuming the insect is large enough to produce the necessary pulling force (if not, the insect may become trapped and die). Pollination is effected by the reverse procedure, in which one of the pollinia becomes trapped within the anther slit. Large-bodied hymenopterans (bees, wasps) are the most common and best pollinators, accounting for over 50% of all Asclepias pollination, whereas monarch butterflies are poor pollinators of milkweed.
Asclepias species produce their seeds in pods termed follicles. The seeds, which are arranged in overlapping rows, bear a cluster of white, silky, filament-like hairs known as the coma (often referred to by other names such as pappus, "floss", "plume", or "silk"). The follicles ripen and split open, and the seeds, each carried by its coma, are blown by the wind. Some, but not all, milkweeds also reproduce by clonal (or vegetative) reproduction.
Selected species
There are also 12 species of Asclepias in South America, among them: A. barjoniifolia, A. boliviensis, A. curassavica, A. mellodora, A. candida, A. flava, and A. pilgeriana.
Ecology
Milkweeds are an important nectar source for native bees, wasps, and other nectar-seeking insects, though non-native honey bees commonly get trapped in the stigmatic slits and die. Milkweeds are also the larval food source for monarch butterflies and their relatives, as well as a variety of other herbivorous insects (including numerous beetles, moths, and true bugs) specialized to feed on the plants despite their chemical defenses.
Milkweeds use three primary defenses to limit damage caused by caterpillars: hairs on the leaves (trichomes), cardenolide toxins, and latex fluids. Data from a DNA study indicate that, generally, more recently evolved milkweed species ("derived" in botany parlance) use these preventive strategies less but grow faster than older species, potentially regrowing faster than caterpillars can consume them.
Research indicates that the very high cardenolide content of Asclepias linaria reduces the impact of the Ophryocystis elektroscirrha (OE) parasite on the monarch butterfly, Danaus plexippus. The OE parasite causes holes to form in the wings of fully developed monarch butterflies. This causes weakened endurance and an inability to migrate. The parasite only infects monarchs when they are larvae and caterpillars, but the detriment is when they are in their butterfly form. By contrast, some species of Asclepias are extremely poor sources of cardenolides, such as Asclepias fascicularis, Asclepias tuberosa, and Asclepias angustifolia.
Monarch butterfly conservation and milkweeds
The leaves of Asclepias species are a food source for monarch butterfly larvae and some other milkweed butterflies. These plants are often used in butterfly gardening and monarch waystations in an effort to help increase the dwindling monarch population.
However, some milkweed species are not suitable for butterfly gardens and monarch waystations. For example, A. curassavica, or tropical milkweed, is often planted as an ornamental in butterfly gardens outside of its native range of Mexico and Central America. Year-round plantings of this species in the United States are controversial and criticised, as they may lead to new overwintering sites along the U.S. Gulf Coast and the consequent year-round breeding of monarchs. This is thought to adversely affect migration patterns, and to cause a dramatic build-up of the dangerous parasite, Ophryocystis elektroscirrha. New research also has shown that monarch larvae reared on tropical milkweed show reduced migratory development (reproductive diapause), and when migratory adults are exposed to tropical milkweed, it stimulates reproductive tissue growth.
Because of this, it is most often suggested to grow milkweeds that are native to the geographical area they are planted in to prevent negative impacts on monarch butterflies.
Monarch caterpillars do not favor butterfly weed (A. tuberosa), perhaps because the leaves of that milkweed species contain very little cardenolide. Some other milkweeds may have similar characteristics.
Uses
Milkweeds are not grown commercially in large scale, but the plants have had many uses throughout human history. Milkweeds have a long history of medicinal, every day, and military use. The Omaha people from Nebraska, the Menomin from Wisconsin and upper Michigan, the Dakota from Minnesota, and the Ponca people from Nebraska, traditionally used common milkweed (A. syriaca) for medicinal purposes.
The bast fibers of some species can be used for rope. The Miwok people of northern California used heart-leaf milkweed (A. cordifolia) for its stems, which they dried and used for cords, strings and ropes.
The fine, silky fluff attached to milkweed seeds, which allows them to be distributed long distances on the wind, is known as floss. Milkweed floss is incredibly difficult to spin due to how short and smooth the filaments are, but blending it with as little as 25% wool or other fiber can produce workable yarn.
A study of the insulative properties of various materials found that milkweed floss was outperformed by other materials in terms of insulation, loft, and lumpiness, but it scored well when mixed with down feathers. The milkweed filaments from the coma (the "floss") are hollow and coated with wax, and have good insulation qualities. During World War II, more than of milkweed floss was collected in the US as a substitute for kapok in life jackets. Milkweed is grown commercially as a hypoallergenic filling for pillows and as insulation for winter coats. Using milkweed floss for these purposes could provide a plant-based alternative to down and promote the growth of milkweed in areas where it has declined, though there is some concern that the environmental impacts could be negative if monoculture is used.
Asclepias is also known as "Silk of America" which is a strand of common milkweed (A. syriaca) gathered mainly in the valley of the Saint Lawrence River in Canada. Milkweed floss can be used in thermal insulation and acoustic insulation. The floss is also highly buoyant and water-repellent, but absorbs oil readily. Due to its oil-absorbing properties, it can be used for oil spill cleanup.
Milkweed latex contains about two percent latex, and during World War II both Nazi Germany and the US attempted to use it as a source of natural rubber, although no record of large-scale success has been found.
Many milkweed species also contain cardiac glycoside poisons that inhibit animal cells from maintaining a proper K+, Ca2+ concentration gradient. As a result, many peoples of South America and Africa used arrows poisoned with these glycosides to fight and hunt more effectively. Some milkweeds are toxic enough to cause death when animals consume large quantities of the plant. Some milkweeds also cause mild dermatitis in some who come in contact with them. Nonetheless, some species can be made edible if properly processed.
| Biology and health sciences | Gentianales | Plants |
480171 | https://en.wikipedia.org/wiki/Common%20myna | Common myna | The common myna or Indian myna (Acridotheres tristis), sometimes spelled mynah, is a bird in the family Sturnidae, native to Asia. An omnivorous open woodland bird with a strong territorial instinct, the common myna has adapted extremely well to urban environments.
The range of the common myna is increasing at such a rapid rate that in 2000 the IUCN Species Survival Commission declared it one of the world's most invasive species and one of only three birds listed among "100 of the World's Worst Invasive Species" that pose a threat to biodiversity, agriculture and human interests. In particular, the species poses a serious threat to the ecosystems of Australia, where it was named "The Most Important Pest/Problem" in 2008.
Taxonomy
In 1760, the French zoologist Mathurin Jacques Brisson included a description of the common myna in his Ornithologie, based on a specimen that he mistakenly believed had been collected in the Philippines. He used the French name Le merle des Philippines and the Latin Merula Philippensis. Although Brisson coined Latin names, they do not conform to the binomial system and are not recognised by the International Commission on Zoological Nomenclature.
When the Swedish naturalist, Carl Linnaeus, updated his Systema Naturae in 1766, for the 12th edition, he added 240 species that had been previously described by Brisson. One of them was the common myna. Linnaeus included a brief description, coined the binomial name Paradisea tristis and cited Brisson's work. The type location was subsequently corrected to Pondicherry in southern India. The specific name tristis is Latin for "sad" or "gloomy". This species is now placed in the genus Acridotheres that was introduced by the French ornithologist Louis Pierre Vieillot in 1816. The generic name Acridotheres is from the Greek ακριδος (akridos), meaning locust, and θηρας (theras), meaning hunter.
Two subspecies are recognised:
the Indian myna (A. t. tristis) (Linnaeus, 1766) – It is found from southern Kazakhstan, Turkmenistan and eastern Iran to southern China, Indochina, the Malay Peninsula and southern India. It has also been introduced to Hawaii and North America. Populations from the northwest of its range have sometimes been separated as a distinct subspecies, A. t. neumanni, while populations from Nepal and Myanmar have been described as A. t. tristoides.
the Sri Lankan myna (A. t. melanosternus) Legge, 1879 – Sri Lanka
The Sri Lankan subspecies melanosternus is darker than the Indian subspecies tristis and has half-black and half-white primary coverts and a larger yellow cheek-patch.
Description
The common myna is readily identified by the brown body, black hooded head and the bare yellow patch behind the eye. The bill and legs are bright yellow. They have rounded wings as well, and a round square tipped tail. There is a white patch on the outer primaries and the wing lining on the underside is white, as well as having a white tail tip. The sexes are similar and birds are usually seen in pairs.
The common myna obeys Gloger's rule in that the birds from northwestern India tend to be paler than their darker counterparts in southern India.
Vocalization
The calls includes croaks, squawks, chirps, clicks, whistles and 'growls', and the bird often fluffs its feathers and bobs its head in singing. The common myna screeches warnings to its mate or other birds in cases of predators in proximity or when it is about to take off flying. Common mynas are popular as cage birds for their singing and "speaking" abilities. Before sleeping in communal roosts, common mynas vocalise in unison, which is known as "communal noise".
Morphometry
Morphometry.
Body length:
Distribution and habitat
The common myna is native to Asia, with its initial home range spanning Iran, Pakistan, India, Nepal, Bhutan, Bangladesh and Sri Lanka, Afghanistan, Uzbekistan, Tajikistan, Turkmenistan, Myanmar, Malaysia, Singapore, peninsular Thailand, Indochina, Japan (both mainland Japan and the Ryukyu Islands) and China.
The common myna has been introduced to many other parts of the world such as Canada, Australia, Israel, New Zealand, New Caledonia, Fiji, the United States (South Florida only), South Africa, Kazakhstan, Kyrgyzstan Uzbekistan, the Cayman Islands, islands in the Indian Ocean (the Seychelles – from which it was subsequently eradicated at great expense, Mauritius, Réunion, Madagascar, the Maldives, the Andaman and Nicobar Islands and the Lakshadweep archipelago) and also in islands of the Atlantic (such as Ascension and Saint Helena, Pacific Ocean and Cyprus February 2022. The range of the common myna is increasing to the extent that in 2000 the IUCN Species Survival Commission declared it among 100 of the world's worst invasive species.
It is typically found in open woodland, cultivation and around habitation. Although it is an adaptable species, its population is abnormal and very much considered a pest in Singapore (where it is locally called as gembala kerbau, literally 'buffalo shepherd') due to competition with the related introduced Javan myna.
The common myna thrives in urban and suburban environments; in Canberra, for instance, 110 common mynas were released between 1968 and 1971. By 1991, common myna population density in Canberra averaged 15 birds per square kilometer. Only three years later, a second study found an average population density of 75 birds per square kilometer in the same area.
The bird likely owes its success in the urban and suburban settings of Sydney and Canberra to its evolutionary origins; having evolved in the open woodlands of India, the common myna is pre-adapted to habitats with tall vertical structures and little to no vegetative ground cover, features characteristic of city streets and urban nature preserves.
The common myna (along with common starlings, house sparrows, and feral rock doves) is a nuisance to city buildings; its nests block gutters and drainpipes, causing water damage to building exteriors.
Behaviour
Breeding
Common mynas are believed to pair for life. They breed through much of the year depending on the location, building their nest in a hole in a tree or wall. They breed at elevations of in the Himalayas.
The normal clutch size is 4–6 eggs. The average size of the egg is . The incubation period is 17 to 18 days and fledging period is 22 to 24 days. The Asian koel is sometimes brood parasitic on this species. Nesting material used by common mynas includes twigs, roots, tow and rubbish. Common mynas have been known to use tissue paper, tin foil and sloughed off snake-skin.
During the breeding season, the daytime activity-time budget of the common myna in Pune in April to June 1978 has been recorded to comprise the following: nesting activity (42%), scanning the environment (28%), locomotion (12%), feeding (4%), vocalisation (7%) and preening-related activities, interactions and other activities (7%).
The common myna uses the nests of woodpeckers, parakeets, etc. and easily takes to nest boxes; it has been recorded evicting the chicks of previously nesting pairs by holding them in the beak and later sometimes not even using the emptied nest boxes. This aggressive behaviour contributes to its success as an invasive species.
There is also some evidence that shows that in introduced environments, the species chooses to nest in more modified and artificial structures than in natural tree cavities when compared to native species.
Food and feeding
Like most starlings, the common myna is omnivorous. It feeds on insects, grubs, earthworms, arachnids, crustaceans, reptiles, small mammals, seeds, grain, fruits, flower nectar and petals, and discarded waste from human habitation. It forages on the ground among grass for insects, and especially for grasshoppers, from which it gets the generic name Acridotheres, "grasshopper hunter". It, however, feeds on a wide range of insects, mostly picked from the ground. It is a cross-pollinator of flowers such as Salmalia and Erythrina. It walks on the ground with occasional hops and is an opportunistic feeder on the insects disturbed by grazing cattle as well as fired grass fields. They prey on eggs and young of other birds, such as Hawaiʻi ʻakepas (Loxops coccineus). They sometimes even wade in shallow waters to catch fish. Living in close proximity to human-made habitats, common mynas may also appear near roadsides to feed on roadkill.
Roosting behaviour
Common mynas roost communally throughout the year, either in pure or mixed flocks with jungle mynas, rosy starlings, house crows, jungle crows, cattle egrets and rose-ringed parakeets and other birds. The roost population can range from less than one hundred to thousands. The time of arrival of mynas at the roost starts before and ends just after sunset. The mynas depart before sunrise. The time and timespan of arrival and departure, time taken for final settlement at the roost, duration of communal sleep, flock size and population vary seasonally.
The function of communal roosting is to synchronise various social activities, avoid predators, exchange information about food sources.
Communal displays (pre-roosting and post-roosting) consist of aerial maneuvers which are exhibited in the pre-breeding season (November to March). It is assumed that this behaviour is related to pair formation.
Invasive species
The IUCN declared the common myna as one of only three birds among the world's 100 worst invasive species (the other two being the red-vented bulbul and the common starling).
The French introduced it in the 18th century from Pondicherry to Mauritius with the aim of controlling insects, even levying a fine on anyone persecuting the bird. It has since been introduced widely elsewhere, including adjacent areas in Southeast Asia, Madagascar, the Middle East, South Africa, the United States, Argentina, Germany, Spain and Portugal, the United Kingdom, Australia, New Zealand and various oceanic islands in the Indian and Pacific Oceans, including prominent populations in Fiji and Hawaii.
The common myna is regarded as a pest in South Africa, North America, the Middle East, Australia, New Zealand and many Pacific islands. It is particularly problematic in Australia. Several methods have been tried to control the bird's numbers and protect native species.
Australia
In Australia, the common myna is an invasive pest. They are often the predominant bird in urban areas along the whole east coast. In a 2008 popular vote, the bird was named "The Most Important Pest/Problem" in Australia. They have earned the nickname "flying rats", due to their numbers and their scavenging behaviour. They are also known as "the cane toad of the sky". However, there is little scientific consensus concerning the extent of its impact on native species.
The common myna was first introduced to Australia between 1863 and 1872, in Victoria, to control insects in the market gardens of Melbourne. At about the same time, the bird is likely to have spread to New South Wales, where it is currently most populous, but documentation is uncertain. The bird was later introduced to Queensland as a predator of grasshoppers and cane beetles. Common myna populations in Australia are now concentrated along the eastern coast around Sydney and its surrounding suburbs, with sparser populations in Victoria and a few isolated communities in Queensland. During 2009, several municipal councils in New South Wales began trials of catching myna birds in an effort to reduce numbers.
The myna can live and breed in a wide range of temperatures, ranging from the frosty winters of Canberra to the tropical climate of Cairns. Self-sustaining populations have been found in regions with a mean monthly highest temperature no less than and a mean monthly lowest temperature no less than , implying that the common myna could spread from Sydney northwards along the eastern coast to Cairns, and westwards along the southern coast to Adelaide, but not to Tasmania, Darwin, or the arid outback regions.
Europe
In 2019, common mynas were added to the List of Invasive Alien Species of Union concern. They have established in Spain and Portugal and were introduced to France, where they occasionally bred.
New Zealand
The common myna was introduced to both the North Island and South Island of New Zealand in the 1870s. However, the cooler summer temperatures in the South Island appear to have impeded the breeding success rate of the southern populations, preventing the proliferation of the species, which was largely non-existent there by the 1890s. In contrast, the North Island population was able to breed more successfully and large portions of the North Island are now populated. However, in the southern reaches of the North Island, the cooler summer temperatures, like those of the South Island, have prevented the establishment of large myna populations. Since the 1950s, mynas have spread northwards and presently inhabit beyond the Waikato region, leading to a majority of its successful population thriving upon lower latitude regions due to the warmer climate. At present, mynas have become especially common in regions of lower latitude, particularly the Northland region, but rarely found south of Whanganui.
South Africa
In South Africa where it escaped into the wild in 1902, it has become very common and its distribution is greater where human populations are greater or where there is more human disturbance. The bird is also notorious for being a pest, kicking other birds out of their nests and killing their young due to the myna's strong territorial instinct. In South Africa it is considered somewhat of a major pest and disturbance of the natural habitat; as a result, it has been declared an invasive species, requiring it to be controlled.
Morphological studies show that the process of spatial sorting is at work on the range expansion of A. tristis in South Africa. Dispersal-relevant traits are significantly correlated with distance from the range core, with strong sexual dimorphism, indicative of sex-biased dispersal. Morphological variations are significant in wing and head traits of females, suggesting females as the primary dispersing sex. In contrast, traits not related to dispersal such as those associated with foraging show no signs of spatial sorting but are significantly affected by environmental variables such as vegetation and intensity of urbanisation.
The United States
In Hawaii, it is out-competing many native birds for food and nesting areas.
To study the invasion genetics and landscape-scale dynamics of A. tristis, scientists have recently developed 16 polymorphic nuclear microsatellite markers using the next generation sequencing (NGS) approach.
Effect on ecosystems and humans
Threat to native birds
The common myna is a hollow-nesting species; that is, it nests and breeds in protected hollows found either naturally in trees or artificially on buildings (for example, recessed windowsills or low eaves). Compared to native hollow-nesting species, the common myna is extremely aggressive, and breeding males will actively defend areas ranging up to 0.83 hectares in size (though males in densely populated urban settings tend to only defend the area immediately surrounding their nests).
This aggressiveness has enabled the common myna to displace many breeding pairs of native hollow-nesters, thereby reducing their reproductive success. In Australia, their aggressiveness has enabled them to chase native birds as large as galahs out of their nests.
The common myna is also known to maintain up to two roosts simultaneously; a temporary summer roost close to a breeding site (where the entire local male community sleeps during the summer, the period of highest aggression), and a permanent all-year roost where the female broods and incubates overnight. Both male and female common mynas will fiercely protect both roosts at all times, leading to further exclusion of native birds.
Threat to crops and pasture
The common myna (which feeds mostly on ground-dwelling insects, tropical fruits such as grapes, plums and some berries and, in urban areas, discarded human food) poses a serious threat to Australian blueberry crops, though its main threat is to native bird species.
In Hawaii, where the common myna was introduced to control pest armyworms and cutworms in sugarcane crops, the bird has helped to spread the robust Lantana camara weed across the islands' open grasslands. It also has been recorded as the fourth-ranking avian pest in the fruit industry by a 2004 survey of the Hawaiian Farm Bureau and the sixth in number of complaints of avian pests overall.
Common mynas can cause considerable damage to ripening fruit, particularly grapes, but also figs, apples, pears, strawberries, blueberries, guava, mangoes and breadfruit. Cereal crops such as maize, wheat and rice are susceptible where they occur near urban areas. Roosting and nesting commensal with humans create aesthetic and health concerns. Common mynas are known to carry avian malaria and exotic parasites such as the Ornithonyssus bursia mite, which can cause dermatitis in humans. The common myna can help spread agricultural weeds: for example, it spreads the seeds of Lantana camara, which has been classed as a Weed of National Significance because of its invasiveness. Common mynas are regularly observed to usurp nests and hollows, destroy the eggs and kill the young of native bird species, including seabirds and parrots. There is evidence that common mynas have killed small land
mammals such as mice, squirrels and possums, but further research on these occurrences is under consideration.
Control
The common myna, being a major agricultural pest and posing a threat to native species in non-native countries, is controlled by various factors. Mynas are either killed or chased away as control. Poison, shooting, cage traps, and bird-scaring devices are currently used for control.
== In culture==
In Sanskrit literature, the common myna has a number of names, most are descriptive of the appearance or behaviour of the bird. In addition to saarika, the names for the common myna include kalahapriya, which means "one who is fond of arguments" referring to the quarrelsome nature of this bird; chitranetra, meaning "picturesque eyes"; peetanetra (one with yellow eyes) and peetapaad (one with yellow legs).
The bird called śārikā () seems to refer to the common myna, though there are other candidates.
Gallery
Explanatory notes
| Biology and health sciences | Passerida | null |
480212 | https://en.wikipedia.org/wiki/Angelshark | Angelshark | Angel sharks are sharks belonging to the genus Squatina. They are the only living members of the family Squatinidae and order Squatiniformes. They commonly inhabit sandy seabeds close to in depth.
Squatina and other Squatiniformes differ from other sharks in having flattened bodies and broad pectoral fins that give them a strong resemblance to rays. They occur worldwide in temperate and tropical seas. Most species inhabit shallow temperate or tropical seas, but a few species inhabit deeper water, down to . Angel sharks are sometimes called monkfish, although this name is also applied to members of the genus Lophius.
While some species occur over a wide geographic range, the majority are restricted to a smaller area. Restriction in geographic range might be as a result of the behaviour of Squatina species, which are ambush predators with a corresponding stationary bottom-dwelling habit. Thus, trans-ocean migration is extremely unlikely, even though large-scale coastal migratory patterns have been reported in species such as Squatina squatina.
Many species are now classified as critically endangered by the International Union for Conservation of Nature. Once common over large areas of the Northeast Atlantic from Norway, Sweden, Morocco and the Canary Islands, to the Mediterranean and Black Seas, fishing pressure has resulted in significant population decline.
Appearance and biology
The angel shark has unique features that differentiates them from other sharks. They are considered as smaller sized sharks because they grow up to only and can weigh around , as opposed to the whale shark that can measure up to and weigh .
While the anterior part of the angel shark's body is broad and flattened, the posterior part retains a muscular appearance more typical of other sharks. The eyes and spiracles are dorsal and the five gill slits are on its back. Both the pectoral and pelvic fins are large and held horizontally. There are two dorsal fins, no anal fin and unusually for sharks, the lower lobe of the caudal fin is longer than the upper lobe. Most types grow to a length of 1.5 m (5 ft), with the Japanese angel shark, known to reach 2 m. Some angel sharks have deformities that have been described in elasmobranchs. These can include skeletal deformities, as lateral spinal curvature (scoliosis), humpback curvature (khyphosis), axial spinal curvature (lordosis), missing fins, additional fins, deformed snout, and more. These abnormalities have only been found in a few sharks, but the causes of these deformities have been found to be from dietary nutritional imbalance, genetic factors, parasites, traumatic injuries, or stress in the specimen. In 2015, two sharks were captured and examined, and both showed a lateral spinal curvature (scoliosis) and also a humpback curvature. Both the animals had the curvature in the middle of their pectoral fins, but the deformity did not affect their swimming capacity.
Spinal scoliosis has been reported to be diverse in sharks, but mostly in pelagic sharks that depend on their swimming abilities to catch their prey. For the angel shark, specifically S. squatina, these curvatures do not seem to significantly affect its hunting capacity, which involves burying itself to ambush their prey. Right now, research is assuming most physical injuries are caused by human interactions because of the constant interference in coastal areas, where most of the sharks reside. There have been few attacks reported, and what few have occurred were due to accidental stepping on of buried newborn sharks. Landings of Pacific angel shark increased through the mid-1980s and reached over 1,125 tonnes in 1986, becoming the shark species with the highest total reported landings off the US West coast that year.
Angel sharks possess extensible jaws that can rapidly snap upwards to capture prey and have long, needle-like teeth. They bury themselves in loose sediment lying in wait for prey, which includes fish, crustaceans and various types of mollusks. They are ovoviviparous, producing litters of up to 13 pups. Pacific angel shark pups are born from March to June in deep water; generally 180 to 300 feet (55 and 90 metres); possibly to protect the pups from predators.
Angel sharks usually reside in depths of and can be seen on muddy or soft benthic substrata where they can easily blend in as they lie in wait. Members of the family Squatinidae have a unique camouflage method, which relates to how they obtain their food, involving lying still on the sea floor, making rapid lunges at passing prey, and using negative pressure to capture prey by sucking it into their mouths.
Species analysis
Morphological identification in the field can be difficult due to discontinuity and similarity of species. In this specific circumstance, the sharks' place within the genus Squatina comprises three species in the southern part of the western Atlantic. The three species observed were Squatina guggenheim, S. occulta and the Brazilian guitarfish Pseudobatos horkelii. These three species are listed in the IUCN Red List as threatened, and they are now protected under Brazilian law, which makes angling and exchange illegal. To prevent landing and trade of these endangered species along the São Paulo, DNA barcoding was used. DNA barcoding revealed fishing and trafficking of these protected species.
Habitat
Angel sharks inhabit temperate and tropical marine environments. They are generally found in shallow waters at depths from off coasts. They are known to bury themselves in sandy or muddy environments during the day, where they remain camouflaged for weeks until a desirable prey crosses paths with them. At night, they take a more active approach and cruise on the bottom of the floor. Squatina preys on fish, crustaceans, and cephalopods.
Behaviour
Although this shark is a bottom-dweller and appears harmless, it can inflict painful lacerations if provoked, due to its powerful jaws and sharp teeth. It may bite if a diver approaches the head or grabs the tail.
Angelsharks have a unique way of breathing compared to most other benthic fish. They do not pump out water from the oropharyngeal cavity like other fish. Instead they use gill flaps located under their body to pump out water during respiration. Doing so also allows them to be more discreet and prevent detection.
Commercial value
Prior to the late 1980s, the Pacific angel shark was considered a "munk fish". It was a byproduct of commercial gillnetting, with no commercial appeal and was used only for crab bait. In 1977, Michael Wagner, a fish processor in Santa Barbara, California, US, in cooperation with local commercial fishermen, developed the market for angel sharks. The annual take of angel shark in 1977 was an estimated 147 kg. By 1985, the annual take of angel shark on the central California coast had increased to more than 454 tonnes or an estimated 90,000 sharks. The population declined dramatically and is now regulated. Angel sharks live very close to shore, resulting in high bycatch rates. In 1991, the use of gillnets in nearshore state waters of California was forbidden, and fishing was restricted in a larger portion of the Pacific angel shark's range.
In April 2008, the UK government afforded the angel shark full protection under the Wildlife and Countryside Act.
Conservation
Once considered abundant in the Atlantic Ocean, the angel shark (Squatina squatina) was classified as "Critically Endangered" in 2010, and recent studies from the IUCN in 2019 reaffirm their CR status. Angel sharks are highly sensitive to bottom trawling and are often caught in gillnets, due to their shallow habitat range.
Angel sharks found in the Mediterranean Sea, S. aculeata, S. oculata, and S. squatina, are at a high risk of extinction, with geographic studies projecting severe population declines for the three species. The Angel Shark Conservation Network, a network established by the IUCN and Shark Trust, is working with authorities from Greece and Turkey to establish conservation strategies to protect angel shark populations in the region.
Evolution
The earliest members of the Squatiniformes are known from the Late Jurassic (from around 160 million years ago) of Europe, assigned to the genus Pseudorhina. Preserved full body specimens of Pseudorhina are very similar to those of living Squatina species. The earliest records that can be assigned with confidence to the modern genus are known from the Early Cretaceous (Aptian) of England.
Species
Currently, the 26 recognized species in this genus are:
Squatina aculeata G. Cuvier, 1829 (sawback angelshark)
Squatina africana Regan, 1908 (African angelshark)
Squatina albipunctata Last & W. T. White, 2008 (eastern angelshark)
Squatina argentina (Marini, 1930) (Argentine angelshark)
Squatina armata (Philippi {Krumweide}, 1887) (Chilean angelshark)
Squatina australis Regan, 1906 (Australian angelshark)
Squatina caillieti J. H. Walsh, Ebert & Compagno, 2011 (Philippines angelshark)
Squatina californica Ayres, 1859 (Pacific angelshark)
Squatina david Acero P, Tavera Vargas, Anguila-Gómez & Hernández-Beracasa, 2016 (David's angelshark)
Squatina dumeril Lesueur, 1818 (sand devil)
Squatina formosa S. C. Shen & W. H. Ting, 1972 (Taiwan angelshark)
Squatina guggenheim Marini, 1936 (angular angelshark)
Squatina heteroptera Castro-Aguirre, Espinoza-Pérez & Huidobro-Campos, 2007 (disparate angelshark)
Squatina japonica Bleeker, 1858 (Japanese angelshark)
Squatina leae Weigmann, Vaz, Akhilesh, Leeney & Naylor 2023 (Lea’s angel shark)
Squatina legnota Last & W. T. White, 2008 (Indonesian angelshark)
Squatina mapama Long, Ebert, Tavera, Acero P., and Robertson, 2021 (Small-crested angelshark)
Squatina mexicana Castro-Aguirre, Espinoza-Pérez & Huidobro-Campos, 2007 (Mexican angelshark)
Squatina nebulosa Regan, 1906 (clouded angelshark)
Squatina occulta Vooren & K. G. da Silva, 1991 (hidden angelshark)
Squatina oculata Bonaparte, 1840 (smoothback angelshark)
Squatina pseudocellata Last & W. T. White, 2008 (western angelshark)
Squatina squatina (Linnaeus, 1758) (angelshark)
Squatina tergocellata McCulloch, 1914 (ornate angelshark)
Squatina tergocellatoides J. S. T. F. Chen, 1963 (ocellated angelshark)
Squatina varii Vaz & Carvalho, 2018 (Brazilian angelshark)
| Biology and health sciences | Sharks | Animals |
480331 | https://en.wikipedia.org/wiki/Siamang | Siamang | The siamang (, ; Symphalangus syndactylus) is an endangered arboreal, black-furred gibbon native to the forests of Indonesia, Malaysia, and Thailand. The largest of the gibbons, the siamang can be twice the size of other gibbons, reaching in height, and weighing up to . It is the only species in the genus Symphalangus. Fossils of siamangs date back to the Middle Pleistocene.
Two features distinguish the siamang from other gibbons. First, two digits on each foot—the second and third toes—are partially joined by a membrane, hence the specific name syndactylus, from the Ancient Greek σύν, sun-, "with" + δάκτυλος, daktulos, "finger". Second, a large gular sac (throat pouch), found in both males and females of the species, can be inflated to the size of the siamang's head, allowing it to make resonating calls.
Two subspecies of the siamang are the nominate Sumatran siamang (S. s. syndactylus) and the Malaysian siamang (S. s. continentis, in Malay Peninsula). Otherwise, the Malaysian individuals are only a population. The siamang occurs sympatrically with other gibbons; its two ranges are entirely within the combined ranges of the agile gibbon and the lar gibbon.
The siamang can live to around 40 years in captivity.
While the illegal pet trade takes a toll on wild populations, the principal threat to the siamang is habitat loss in both Indonesia and Malaysia. The palm oil production industry is clearing large swaths of forest, reducing the habitat of the siamang, along with those of other species, such as the Sumatran tiger.
Description
The siamang has long, dense, shaggy hair, which is the darkest shade of all gibbons. The ape's long, gangling arms are longer than its legs. The average length of a siamang is 90 cm; the largest they have ever grown is 150 cm. The face of this large gibbon is mostly hairless, apart from a thin mustache.
Distribution and habitat
The siamang inhabits the forest remnants of Sumatra Island and the Malay Peninsula, and is widely distributed from lowland forest to mountain forest—even rainforest—and can be found at altitudes up to 3800 m. It lives in groups of up to six individuals (four individuals on average) with an average home range of 23 hectares. Their day ranges are substantially smaller than those of sympatric Hylobates species, often less than 1 km. The siamang's melodious singing breaks the forest's silence in the early morning after the agile gibbons' or lar gibbons' calls. The siamangs in Sumatra and the Malay Peninsula are similar in appearance, but some behaviors differ between the two populations.
Ecology and behavior
Siamangs have an ecology and relationship between two types of gibbons that share the same habitat. Those include the Agile gibbon and Lar gibbon. Both of the gibbons that live with the siamangs are Hylobates rather than Symphalangus. When two siamangs meet, they often have a bond with each other. They might also communicate by using their throat pouches and shouts to communicate when they feel excited, relaxed, trying to mate, or threatened. If siamangs use loud sounds, they also use body language to communicate. They use sign language or pointing to make others aware of what they need or what they want to do.
Diet
The siamang eats mainly various parts of plants. The Sumatran siamang is more frugivorous than its Malayan relative, with fruit making up to 60% of its diet. The siamang eats at least 160 species of plants, from vines to woody plants. Its major food source is figs (Ficus spp.). The siamang prefers to eat ripe rather than unripe fruit, and young rather than old leaves. It eats flowers and a few animals, mostly insects. When the siamang eats large flowers, it eats only the corollae (petals), but it eats all parts of smaller flowers, with the small fruit collected in its hand before being consumed. When it eats big and hard seeds or seeds with sharp edges, it peels out the fruit flesh and throws away the seed. Although its diet consists of substantial portions of fruit, it is the most folivorous of all members of Hylobatidae. As it is also the largest gibbon, it fits well with the general primate dietary trend in which larger primates tend to be more folivorous.
Demography and population
A group of siamangs normally consists of an adult dominant male, an adult dominant female, with offspring, infants, and sometimes a subadult. The subadult usually leaves the group after attaining the age of 6–8 years; subadult females tend to leave the group earlier than subadult males. Siamang gestation period is between 6.2 and 7.9 months; after the infant is born, the mother takes care of the infant for the first year of its life. Siamang males tend to offer more paternal care than do other members of the family Hylobatidae, taking up a major role in carrying an infant after it is about 8 months old. The infant typically returns to its mother to sleep and nurse. The infant begins to travel independently from its parents by its third year of life.
Siamangs are generally known to have monogamous mating pairs, which have been documented to spend more time in close proximity to each other, in comparison to other gibbon species. Both monogamous and polyandrous groups, though, are found in South Sumatra. In studying these populations, infants belonging to monogamous groups were found to receive more overall male care than infants in the polyandrous groups. This reduced care is most likely due to reduced certainty of paternity in these groups.
Habitat disturbance affects siamang group composition; it is varied in age-sex structure between intact forest and burnt, regrown forest. The burnt, regrown forest population contained more adult and subadults than the intact forest population, which had more infants, small juveniles, and large juveniles. Infant survival rates in burnt, regrown forest groups are lower than in intact forest groups. The number of individuals in the latter is higher than in the former. The siamang in disturbed forests live in small groups and have a density lower than in intact forests because of lack of food resources and trees for living.
In the 1980s, the Indonesian population of the siamang in the wild was estimated to be 360,000 individuals. This figure may be less in the 21st century: Bukit Barisan Selatan National Park is the third-largest protected area () in Sumatra, of which roughly remain under forest cover inhabited by 22,390 siamangs (in 2002 censuses). In Sumatra, the siamang prefers to inhabit lowland forest between above sea level.
Behavior
The siamang tends to rest for more than half of its waking period from dawn to dusk, followed by feeding, moving, foraging, and social activities. It takes more rest during midday, taking time to groom others or to play. During resting time, it usually uses a branch of a large tree, lying on its back or belly. Feeding behaviors, foraging, and moving are most often in the morning and after resting. Grooming is one of the most important social interactions among family members. Grooming takes place between adults earlier in the day; the adults groom the juveniles later in the day. Adult males are the most involved in grooming.
In the dry season, the size of the siamang's daily range is larger than in the rainy season. The siamang in southern Sumatra spends less time foraging than siamangs in other places, as it eats a diet higher in fruit. It thus consumes more nutrients, which results in less time needed for food acquisition. A siamang may spend an entire day in a single fruiting tree, moving out when it wants to rest and returning to feed.
Siamangs are a very social species of primates and exhibit a variety of tactile and visual gestures, along with actions and facial expressions to communicate and increase social bonds within their family group. They are also territorial, and interact with other family groups by making loud calls to let other groups know where their territory is. The calls may be asynchronous, where they are not directed at a particular neighbouring group, or simultaneous group calls may take place across the territory boundary. Males are known to chase one another across the boundary.
Grooming frequency between males and females has been found to correlate to copulation frequency, as well as bouts of aggression. Pairs copulate over four to five months at intervals of two to three years. The peak of their reproductive activity is often during the time when fruit is most abundant. Dorsoventral copulation is the most common type in siamangs, where the female is squatting and the male hangs by his arms and grips the female with his legs, whereas ventroventral copulation, where both primates are suspended, occurs only one in 60 times on average.
Role of calling
The siamang starts its day by calling in the early morning; it calls less after midday, with the peak of the calls around 9:00 to 10:00 am. Most of the siamang's calls are directed to its neighbours rather than to those inside its home range. This means the siamang's calling is in response to disturbances and to defend its territory. Calls in the late morning typically happen when it meets or sees another siamang group. The edge of the siamang's home range, which may overlap another, is often the place where calling is made. Counter (co-response) calling occasionally happens near the border or in the overlap area. Calls are numerous when fruit is more abundant rather than when it is less available. Branch shaking, swinging, and moving around the tree crowns accompany the calling. This movement might be to show the other groups where they are.
The siamang prefers calling in the living, tall, and big trees, possibly where another group is easy to see. Besides that, such trees can support siamang movement. Calling trees are usually near feeding trees, but sometimes they call in the feeding trees.
Mated pairs produce loud, well-patterned calling bouts, which are referred to as duetting. These calls advertise the presence and status of a mated pair. Newly formed pairs spend more time singing than an established pair. Advertising the presence of a strong bond is advantageous in territorial defense. Siamang duetting differs from other species because it has a particularly complex vocal structure. Four distinct classes of vocalizations have been documented: booms, barks, ululating screams, and bitonal screams. Females typically produce long barks and males generally produce bitonal screams, but both sexes have been known to produce all four classes of vocalizations. Unlike other gibbons in which vocalization is added by laterally expanded laryngeal sacs, those in siamangs fuse with each other and extend into the ventral area of the neck.
Seeding
As a frugivorous animal, the siamang disperses seeds through defecation as it travels across its territory. The siamang can carry seed while digesting, and defecate between from the seed resource, which supports the forest's regeneration and succession.
Threats and conservation
As an arboreal primate whose survival absolutely depends on the forest, the siamang faces population pressure due to habitat loss, poaching, and hunting.
Habitat loss
A major threat to the siamang is habitat fragmentation due to plantations, forest fire, illegal logging, encroachment, and human development. Firstly, palm-oil plantations have removed large areas of the siamang's habitat in recent decades. Since 2002, 107,000 km2 of oil palm have been planted, which has replaced much rainforest in Indonesia and Malaysia, where the siamang originally lived. Secondly, in the second decade of the 21st century, forests in the Malay Peninsula have been destroyed due to illegal logging. Sixteen out of the 37 permanent forest reserves in Kelantan, in the Malay Peninsula, where most of the siamangs live, have been encroached upon by illegal loggers. Thirdly, forest encroachments change forest cover into cultivated land; for example, the rising price of coffee in 1998 encouraged people in Sumatra to replace the forest with coffee plantations. Fourthly, development in many areas requires infrastructure, such as roads, which now divide conservation areas and have caused forest fragmentation and edge effects.
Poaching and hunting
Unlike other parts of Asia, primates are not hunted for their meat in Indonesia. They are hunted for the illegal pet trade, with hunters preferring infant siamangs. Poachers often kill the mothers first, since siamang females are highly protective of their infants, and removing the infant without first killing the mother requires more effort. Most siamangs on the market are infants, which often die during transportation.
Conservation
Siamang can be found in at least 11 protected areas:
Indonesia
Bukit Barisan Selatan National Park
Gunung Leuser National Park
Kerinci Seblat National Park
Langkat Barat Wildlife Reserve
Way Kambas National Park
Malaysia
Fraser's Hill Reserve
Gunong Besout Forest Reserve
Krau Wildlife Reserve
Ulu Gombak Wildlife Reserve
Thailand
Hala Bala Wildlife Sanctuary
| Biology and health sciences | Apes | Animals |
480516 | https://en.wikipedia.org/wiki/Aspleniaceae | Aspleniaceae | The Aspleniaceae (spleenworts) are a family of ferns, included in the order Polypodiales. The composition and classification of the family have been subject to considerable changes. In particular, there is a narrow circumscription, Aspleniaceae s.s. (adopted here), in which the family contains only two genera, and a very broad one, Aspleniaceae s.l., in which the family includes 10 other families kept separate in the narrow circumscription, with the Aspleniaceae s.s. being reduced to the subfamily Asplenioideae. The family has a worldwide distribution, with many species in both temperate and tropical areas. Elongated unpaired sori are an important characteristic of most members of the family.
Description
Members of the family grow from rhizomes, that are either creeping or somewhat erect, and are usually but not always unbranched, and have scales that usually have a lattice-like (clathrate) structure. In some species, for example Asplenium nidus, the rhizomes form a kind of basket which collects detritus. The leaves may be undivided or be divided, with up to four-fold pinnation. The sori are characteristic of the family. They are elongated, and normally located on one side of a vein. More rarely, they may be in pairs on a single vein, but then they never curve over the vein. A flap-like indusium arises along one edge of a sorus. The leaf stalks (petioles) have two vascular bundles, uniting to form an X-shape in cross-section towards the tip of the leaf. The stalks of the sporangia are one cell wide in the middle.
Taxonomy
The family Aspleniaceae was first described by Edward Newman in 1840. Newman included three genera: Athyrium, Asplenium and Scolopendrium. Athyrium is now placed in a different family, Athyriaceae, not considered very strongly related to the Aspleniaceae, and Scolopendrium is regarded as synonym of Asplenium.
The narrow circumscription of the family adopted by the Pteridophyte Phylogeny Group classification of 2016 (PPG I) recognizes only two genera, Asplenium and Hymenasplenium. Asplenium has previously been split into a dozen or so genera, including Diella, found only in Hawaii. The consensus of molecular phylogenetic studies is that all are nested within Asplenium. PPG I places Aspleniaceae in the suborder Aspleniineae of the order Polypodiales.
Earlier, Christenhusz and Chase had proposed a much broader circumscription of Aspleniaceae, in which it consisted of all the separate families that PPG I places in the suborder Aspleniineae (eight at the time), with the families reduced to subfamilies. Thus the Aspleniaceae of PPG I became the subfamily Asplenioideae. , the broader circumscription of the Aspleniaceae is used by Plants of the World Online, which lists 24 genera.
Phylogenic relationships
Aspleniaceae is placed in a clade known as eupolypods II, or more formally as suborder Aspleniineae. The following cladogram, based on Lehtonen (2011) and Rothfels & al. (2012), shows a likely phylogenic relationship between the Aspleniaceae and the other families in the clade.
Genera
In the PPG I system, Aspleniaceae s.s. contains two genera:
Asplenium L. – about 700 species, worldwide
Hymenasplenium Hayata – about 40 species, tropical and subtropical
Distribution and habitat
The Aspleniaceae have a worldwide distribution, with the large genus Asplenium being native to almost all parts of the world except Antarctica and some high Arctic areas. The family is unusual in having high diversity in both temperate and tropical areas, and more-or-less equal numbers of terrestrial and epiphytic species. Plants are terrestrial, growing in the ground, lithophytic, growing on rocks, or epiphytic, growing on other plants; less often they are aquatic, growing in moving water.
| Biology and health sciences | Ferns | Plants |
480586 | https://en.wikipedia.org/wiki/Tarpan | Tarpan | The tarpan (Equus ferus ferus) was a free-ranging horse population of the Eurasian steppe from the 18th to the 20th century. What qualifies as a tarpan is subject to confusion. It is unknown whether those horses represented genuine wild horses, feral domestic horses or hybrids. The last individual believed to be a tarpan died in captivity in the Russian Empire in 1909.
Beginning in the 1930s, several attempts were made to develop horses that looked like tarpans through selective breeding, called "breeding back" by advocates. The breeds that resulted included the Heck horse, the , and a derivation of the Konik breed, all of which have a primitive appearance, particularly in having the grullo coat colour. Some of these horses are now commercially promoted as "tarpans", although such animals are only domestic breeds and not the wild animal themselves.
Name and etymology
The name "tarpan" or "tarpani" derives from a Turkic language (Kazakh or Kyrgyz) name meaning "wild horse". The Tatars and the Cossacks distinguished the wild horse from the feral horse; the latter was called Takja or Muzin.
Taxonomy
The tarpan was first described by Samuel Gottlieb Gmelin in 1771; he had seen the animals in 1769 in the district of Bobrov, near Voronezh. In 1784, Pieter Boddaert named the species Equus ferus, referring to Gmelin's description. Unaware of Boddaert's name, Otto Antonius published the name Equus gmelini in 1912, again referring to Gmelin's description. Since Antonius' name refers to the same description as Boddaert's it is a junior objective synonym. It is now thought that the domesticated horse, named Equus caballus by Carl Linnaeus in 1758, is descended from the tarpan; indeed, many taxonomists consider them to belong to the same species. By a strict application of the rules of the International Code of Zoological Nomenclature, the tarpan ought to be named E. caballus, or if considered a subspecies, E. caballus ferus. However, biologists have generally ignored the letter of the rule and used E. ferus for the tarpan to avoid confusion with the domesticated subspecies.
It is debated if the small, free-roaming horses seen in the Russian steppes during 18th and 19th centuries and called "tarpan" were indeed wild, never-domesticated horses, hybrids of the Przewalski's horse and local domestic animals, or simply feral horses. Most studies have been based on only two preserved specimens and research to date has not positively linked the tarpan to Pleistocene or Holocene-era animals.
In 2003, the International Commission on Zoological Nomenclature "conserved the usage of 17 specific names based on wild species, which are predated by or contemporary with those based on domestic forms", confirming E. ferus for the undomesticated wild horse. Taxonomists who consider the domestic horse a subspecies of the wild horse should use Equus ferus caballus; the name Equus caballus remains available for the domestic horse where it is considered to be a separate species.
Appearance
Traditionally, two tarpan subtypes have been proposed, the forest tarpan and steppe tarpan, although there seem to be only minor differences in type. The general view is that there was only one subspecies, the tarpan, Equus ferus ferus. The last individual, which died in captivity in 1909, was between tall at the shoulders, and had a thick, falling mane, a grullo coat colour, dark legs, and primitive markings, including a dorsal stripe and shoulder stripes.
A number of coat colour genotypes have been identified within European wild horses from the Pleistocene and Holocene: those creating bay, black and leopard complex are known from the wild horse population in Europe and were depicted in cave paintings of wild horses during the Pleistocene. The dun gene, a dilution gene seen in Przewalski's horse that also creates the grullo or "blue dun" coat seen in the Konik, has not yet been studied in European wild horses. It is likely that at least some wild horses had a dun coat.
Historical reports are ambiguous on whether tarpans had standing manes like wild equines, or falling manes like domestic horses.
The appearance of European wild horses may be reconstructed with genetic, osteologic and historic data. One genetic study suggests that bay was the predominant colour in European wild horses. During the Mesolithic, a gene coding a black coat colour appeared on the Iberian peninsula. This colour spread east but was less common than bay in the investigated sample and never reached Siberia. Bay in combination with dun results in the "bay dun" colour seen in Przewalski's horses; black with dun creates the grullo coat. A loss of the dun dilution may have been advantageous in more forested western European landscapes, as dark colours were a better camouflage in forests. Pangaré or "mealy" coloration, a characteristic of other wild equines, might have been present in at least some European wild horses, as historic accounts report a light belly.
Historic references report that most tarpans were black dun (grullo, "blue dun") colour. Some black individuals were reported to have domestic colours like white or grey legs. Authors, such as Peter Pallas, believed tarpans to be escaped farm horses.
History
Wild horses have been present in Europe since the Pleistocene and ranged from southern France and Spain east to central Russia. There are cave drawings of primitive predomestication horses at Lascaux, France and in Cave of Altamira, Spain, as well as artifacts believed to show the species in southern Russian empire, where a horse of this type was domesticated around 3000 BC. Equus ferus had a continuous range from western Europe to Alaska; historic material suggests wild horses lived in most parts of Holocene continental Europe and the Eurasian steppe, except for parts of Scandinavia, Iceland and Ireland.
DNA
DNA obtained from a "tarpan" that lived in the steppes of the Kherson region during the mid-19th century shows its ancestry to be a mixture between horses native to Europe and found with the Corded Ware Culture, and horses closely related to the DOM2 population, the ancestors of modern domestic horses. (DOM2 horses originated from the Western Eurasia steppes, especially the lower Volga-Don, but not in Anatolia, during the late fourth and early third millennia BCE. Their genes may show selection for easier domestication and stronger backs). This is not consistent with tarpans as the wild ancestor of, or a feral version of, DOM2. Nor is the tarpan a hybrid with Przewalski's horses.
Forest tarpan
The "forest horse" or "forest tarpan" was a hypothesis of various 20th-century natural scientists, including Tadeusz Vetulani, who suggested that the continuous forestation of Europe after the last ice age created forest-adapted subtype of the wild horse, which he named Equus sylvestris. However, historic references do not describe any major difference between the populations, and therefore most authors assume there was only one subspecies of western Eurasian wild horse, Equus ferus ferus.
Nevertheless, a stocky type of horse living in forests and highlands was described during the 19th century in Spain, the Pyrenees, the Camargue, the Ardennes, Great Britain, and the southern Swedish upland.
They had a robust head and strong body, and a long frizzy mane. The colour was described as faint brown or yellowish brown with eel stripe and leg stripes, or wholly black legs. The flanks and shoulders were spotted, some of them tended to an ashy colour. They dwelled in rocky habitats and showed intelligent and fierce behaviour. Yet, those horses were never colloquially called "tarpans".
Black wild horses were found in Dutch swamps, with a large skull, small eyes, and a bristly muzzle. Their mane was full, with broad hooves, and curly hair. However, it is possible that these were feral and not wild horses.
Herodotus described light-coloured wild horses in an area now part of Ukraine in the 5th century BCE. In the 12th century, Albertus Magnus stated that mouse-coloured wild horses with a dark eel stripe lived in the German territory, and in Denmark, large herds were hunted.
16th century free-ranging horses in Europe
Wild horses still were common in the east of Prussia during the 15th and early 16th centuries. During the 16th century, wild horses disappeared from most of the mainland of western Europe and became less common in eastern Europe as well. Belsazar Hacquet saw wild horses in the Polish zoo in Zamość during the Seven Years' War. According to him, those wild horses were of small body size, had a blackish brown colour, a large and thick head, short dark manes and tail hair, and a "beard". They were absolutely untameable and defended themselves harshly against predators.
Kajetan Kozmian visited the population at Zamość as well and reported that they were small and strong, had robust limbs and a consistently dark mouse colour. Samuel Gottlieb Gmelin witnessed herds in Voronezh in 1768. Those horses were described as very fast and shy, fleeing at any noise, and as small, with small, pinned ears and a short frizzly mane. The tail was shorter than in domestic horses. They were typically mouse-coloured with a light belly and legs becoming black, although grey and white horses were mentioned as well. The coat was long and dense. The horses at Zamosc were never called "tarpan" back in their lifetime.
Peter Simon Pallas witnessed possible tarpans in the same year in southern Russia. He thought they were feral animals that escaped during the confusions of wars. These herds were important game of the Tatars and numbered between 5 and 20 animals. The horses he described had a small body, large and thick heads, short frizzly manes and short tail hair, as well as pinned ears. The colour was described as faint brownish, sometimes brown or black. He also reported of obvious domestic hybrids with light-coloured legs or grey coats.
The tarpans of 18th century Europe
The Natural History of Horses by 19th-century author Charles Hamilton Smith also described tarpans. According to Smith, the herds of free-ranging horses numbered from a few to hundred individuals. They often were mixed with domestic horses, and alongside pure herds there were herds of feral horses or hybrids. The colour of alleged pure tarpans was described as consistently brown, cream-coloured or mouse-coloured. The short frizzy mane was reported to be black, as were the tail and legs. The ears were of varying size, but set high on the skull. The eyes were small.
According to Smith, tarpans made stronger sounds than domestic horses and the overall appearance of these horses was mule-like.
Extinction
The human-caused extinction of wild horses in Europe began in Southern Europe, possibly in antiquity. While humans had been hunting wild horses since the Paleolithic, during historic times horse meat was an important source of protein for many cultures. As large herbivores, the range of the tarpan was continuously decreased by the increasing human population of the Eurasian landmass. Wild horses were further targeted because they caused damage to hay stores and often took domestic mares from pastures. Furthermore, interbreeding with wild horses was an economic loss for farmers since the foals of such matings were intractable.
Tarpans lived in the southern parts of the steppe. In 1879 the last scientifically confirmed individual was killed. After that, only dubious sightings were documented. As the tarpan horse died out in the wild between 1875 and 1890, the last free-ranging mare considered a tarpan was accidentally killed during an attempt at capture. The last captive tarpan died in 1909 in a Ukrainian zoo in the Russian empire.
The horses at Zamosc survived until 1806, when they were allegedly donated to local farmers of Poland and bred to their horses. The Konik is claimed to descend from these hybrid horses, as recent research has highlighted a significant degree of anatomic difference between free-roaming Konik in the Netherlands and other modern domesticated horses.
Tarpans interbreeding with domestic horses
The oldest archaeological evidence for domesticated horses is from Kazakhstan and Ukraine between 6,000 and 5,500 YBP (years before present). The diverse mitochondrial DNA of domestic horses contrasts sharply with the very low diversity of the Y chromosome; that suggests that many mares but only a few stallions were used, and local use of wild mares or even secondary sites of domestication are likely. Therefore, the European wild horse may well have contributed to the domestic horse.
Wild horses vs. feral horses
Some researchers consider the tarpans to be mixed wild and feral population or completely feral horses. Few consider the more recent animals historically called "tarpans" to be genuine wild horses without domestic influence. Historic references to "wild horses" may actually refer to feral domestic horses or hybrids.
Some 19th century authors wrote that local "wild" horses had hoof problems that led to crippled legs; therefore, they assumed these were feral horses. Other contemporary authors claimed all "wild" horses between the Volga River and the Ural were actually feral. However, others thought that this was too speculative and assumed that wild, undomesticated horses still lived into the 19th century. Domestic horses used in warfare often were turned loose when they were no longer needed. Also, remaining wild stallions could steal domestic mares. There are some accounts from the 18th and 19th centuries of wild herds with typical wild horse features such as large heads, pinned ears, short frizzy mane and tail, but mentioned animals with domestic influence as well.
The only known individual to be photographed was the so-called Cherson tarpan, which was caught as a foal near Novovorontsovka in 1866. It died in 1887 in the Moscow Zoo. The nature of this horse was dubious in its lifetime, because it showed almost none of the wild horse features described in the historic sources. In 2021, a study found that the so-called 'Shatilov' tarpan, a museum specimen from the Kherson region of the Pontic–Caspian steppe that died in 1868, was a hybrid between two horse lineages, with two thirds of its genetics representing the same ancestral lineage as modern horses, and the remaining third related to a distinct lineage of European horses found at an archaeological site associated with the late-Neolithic Corded Ware culture.
Breeding back the tarpan
Three attempts have been made to use selective breeding to create a type of horse that resembles the tarpan phenotype, though recreating an extinct subspecies is not genetically possible with current technology. In 1936, Polish university professor Tadeusz Vetulani selected Polish farm horses that were formerly known as Panje horses (now called Konik) and that he believed resembled the historic tarpan and started a selective breeding program. Although his breeding attempt is well known, it made only a minor contribution to the modern Konik stock, which clusters genetically with other domestic horse breeds, including those as diverse as the Mongolian horse and the Thoroughbred.
In the early 1930s, Berlin Zoo Director Lutz Heck and Heinz Heck of the Munich Zoo began a program crossbreeding Koniks with Przewalski horse stallions, and the mares of Gotland ponies, and Icelandic horses,. By the 1960s they produced the Heck horse. In the mid-1960s, Harry Hegard started a similar program in the United States using mustangs and local working ranch horses that has resulted in the Hegardt or Stroebel's horse.
Assessment
While all three breeds have a primitive look that resembles the tarpan in some respects, they are not genetically tarpans and the wild, predomestic European horse remains extinct. However, this does not prevent some modern breeders from marketing horses with these features as a "tarpan".
In spite of sharing primitive external features, the Konik and Hucul horses have markedly different conformation with differently proportioned body measurements, thought in part to be linked to living in different habitats.
Other breeds sometimes alleged to be surviving wild horses include the Exmoor pony and the Dülmen pony. However, genetic studies do not set any of these breeds apart from other domestic horses. On the other hand, there has not yet been a study comparing domestic breeds directly with the European wild horse.
| Biology and health sciences | Equidae | Animals |
480641 | https://en.wikipedia.org/wiki/Glanders | Glanders | Glanders is a contagious zoonotic infectious disease that occurs primarily in horses, mules, and donkeys. It can be contracted by other animals, such as dogs, cats, pigs, goats, and humans. It is caused by infection with the bacterium Burkholderia mallei.
Glanders is endemic in Africa, Asia, the Middle East, and Central and South America. It has been eradicated from North America, Australia, and most of Europe through surveillance and destruction of affected animals and import restrictions. It has not been reported in the United States since 1945, except in 2000, when an American lab researcher had an accidental exposure in the lab. It is a notifiable disease in the UK, although it has not been reported there since 1928.
The term is from Middle English or Old French , both meaning glands. Other terms include , , and .
Presentation
Signs of glanders include the formation of nodular lesions in the lungs and ulceration of the mucous membranes in the upper respiratory tract. The acute form results in coughing, fever, and the release of an infectious nasal discharge, followed by septicaemia and death within days. In the chronic form, nasal and subcutaneous nodules develop, eventually ulcerating; death can occur within months, while survivors act as carriers.
Cause and transmission
Glanders is caused by infection with the Burkholderia mallei, usually by ingestion of contaminated feed or water. B. mallei is able to infect humans, so it is classed as a zoonotic agent. Transmission occurs by direct contact with infected animal's body fluid and tissues and entry is through skin abrasions, nasal and oral mucosal surfaces, or inhalation.
Diagnosis
The mallein test is a sensitive and specific clinical test for glanders. Mallein (ATCvet code: ), a protein fraction of the glanders organism (B. mallei), is injected intradermopalpebrally or given by eye drop. In infected animals, the eyelid swells markedly in 1 to 2 days.
Historical cases and potential use in war
Glanders has been known since antiquity, with a description by Hippocrates around 425 BCE.
From the Middle Ages to the 1900s, glanders was a significant threat to armies. Before the Battle of Blenheim in 1704, glanders may have afflicted and greatly diminished the horses of Marshal Tallard's cavalry, helping the Duke of Marlborough win the battle.
Glanders was a significant problem for civilian use of horses, as well. In the 18th-century veterinary hospital at the École Nationale Vétérinaire d'Alfort, glanders was the most common disease among their equine patients and the one most likely to cause death.
Due to the high mortality rate in humans and the small number of organisms required to establish infection, B. mallei is regarded as a potential biological warfare or bioterrorism agent, as is the closely related organism, B. pseudomallei, the causative agent of melioidosis. During World War I, glanders was believed to have been spread deliberately by German agents to infect large numbers of Russian horses and mules on the Eastern Front. Other agents attempted to introduce the disease in the United States and Argentina. This had an effect on troop and supply convoys, as well as on artillery movement, which were dependent on horses and mules. Human cases in Russia increased with the infections during and after WWI. The Japanese deliberately infected horses, civilians, and prisoners of war with B. mallei at the Unit 731 Pingfang (China) Institute and Unit 100 facilities during World War II. The U.S. studied this agent as a possible biological weapon in 1943–44, but did not weaponize it. U.S. interest in glanders (agent LA) continued through the 1950s, except it had an inexplicable tendency to lose virulence in the lab, making it difficult to weaponize. Between 1982 and 1984, the Soviet Union allegedly used weaponized B. mallei during the Soviet–Afghan War.
Vaccine research
No vaccine is licensed for use in the U.S. Infection with these bacteria results in nonspecific symptoms and can be either acute or chronic, impeding rapid diagnosis. The lack of a vaccine for either bacterium also makes them potential candidates for bioweaponization. Together, with their high rate of infectivity by aerosols and resistance to many common antibiotics, both bacteria have been classified as category B priority pathogens by the US NIH and US CDC, which has spurred a dramatic increase in interest in these microorganisms. Attempts have been made to develop vaccines for these infections, which would not only benefit military personnel, a group most likely to be targeted in an intentional release, but also individuals who may come in contact with glanders-infected animals or live in areas where melioidosis is endemic.
| Biology and health sciences | Bacterial infections | Health |
4810606 | https://en.wikipedia.org/wiki/Ulmus%20minor | Ulmus minor | Ulmus minor Mill., the field elm, is by far the most polymorphic of the European species, although its taxonomy remains a matter of contention. Its natural range is predominantly south European, extending to Asia Minor and Iran; its northern outposts are the Baltic islands of Öland and Gotland, although it may have been introduced by humans. The tree's typical habitat is low-lying forest along the main rivers, growing in association with oak and ash, where it tolerates summer floods as well as droughts.
Current treatment of the species owes much to Richens, who noted (1983) that several varieties of field elm are distinguishable on the European mainland. Of these, he listed the small-leaved U. minor of France and Spain; the narrow-leaved U. minor of northern and central Italy; the densely hairy leaved U. minor of southern Italy and Greece; the U. minor with small-toothed leaves from the Balkans; the U. minor with large-toothed leaves from the Danube region; and the small-leaved U. minor from southern Russia and Ukraine. As for British varieties, "the continental populations most closely related [to eastern English Field Elm] are in central Europe", while south-western forms were introduced from France. He concluded, however, that owing to incomplete field-research at the time of writing, it was "not possible to present an overall breakdown of the European Field Elm into regional varieties". The epithet 'red' elm was commonly used by British foresters, an allusion to the colour of the timber.
Richens sank a number of British elms, notably English elm, as either subspecies or varieties of U. minor in 1968. However, Melville, writing ten years later, identified five distinct species (including U. glabra in the count), several varieties and numerous complex hybrids. In 1992 Armstrong identified no fewer than forty British species and microspecies. Clive Stace (1997) wrote of the British elms "The two-species (glabra and minor) concept of Richens is not sufficiently discriminating to be of taxonomic value". Nevertheless, it is Richens’ classification which has been the most commonly adopted in recent years, although it is not used in Flora Europaea.
In 2009 Dr Max Coleman of the Royal Botanic Garden, Edinburgh wrote: "The advent of DNA fingerprinting has shed considerable light on the question. A number of studies have now shown that the distinctive forms Melville elevated to species and Richens lumped together as field elm are single clones, all genetically identical, which have been propagated by vegetative means such as cuttings or root suckers. This means that enigmatic British elms such as Plot elm and English elm have been shown to be single clones of field elm. Although Richens did not have the evidence to prove it, he was correct in recognising a series of clones and grouping them together as a variable species."
It is hoped that analysis of molecular markers will ultimately eliminate the taxonomic confusion.
Description
The tree typically grows to < 30 m (98 ft) and bears a rounded crown. The bark of the trunk is rough, furrowed lightly in older trees to form a block pattern. Young branchlets occasionally have corky wings. The shoots are slender compared with those of wych elm. The leaves are smaller than those of the other European species, hence the specific epithet minor, however they can vary greatly according to the maturity of the tree. Leaves on juvenile growth (suckers, seedlings etc.) are coarse and pubescent, whereas those on mature growth are generally smooth, though remaining highly variable in form; there are generally fewer than 12 pairs of side veins. A common characteristic is the presence of minute black glands along the leaf veins, detectable with the aid of a magnifying glass. The samarae are oval or obovate, glabrous, long, notched at the top, with the seed close to the notch. Ulmus minor in France generally begins to flower and fruit when aged 10 years.
The species readily produces suckers from roots and stumps, even after devastation by Dutch elm disease; consequently genetic resources are not considered endangered.
Pests and diseases
The species has a hugely variable reaction to Dutch elm disease (DED), including all the fashionable pre-20th century plantsman's clones (see Subspecies and varieties). However, field elm is genetically highly variable; Italian specimens when inoculated with the pathogen displayed between 15 and 100% dieback and between 70 and 100% wilting, whereas with trees tested in Spain, the variability ranged from 5 to 100% dieback, and 20 to 95% wilting. In 2013 researchers at the Universidad Politėcnica de Madrid announced the discovery and cloning of trees in Spain with levels of resistance greater than 'Sapporo Autumn Gold' (see Cultivation).
Tolerance of elm yellows (phloem necrosis) is generally good, U. minor exhibiting symptoms such as the 'witch's broom' only sporadically throughout Italy, including Sicily and Sardinia, however the disease was often locally common within the species in France, including Paris.
Cultivation
U. minor in general and a number of clones in particular (see 'Cultivars' below) were once commonly cultivated across Europe in town and country, but owing to its susceptibility to Dutch elm disease, U. minor is now uncommon in cultivation. However, in an ongoing project that began in the 1990s, several thousand surviving field elms have been tested for innate resistance by national research institutes in the EU, with a view to returning field elm to cultivation. Results from Spain (2013), for example, confirm that a very small number of surviving field elms (about 0.5% of those tested) appear to have comparatively high levels of tolerance of the disease, and it is hoped that a controlled crossing of the best of these will produce resistant Ulmus minor hybrids for cultivation.
In the UK, despite its late leaf-flush in the north and its suckering habits, continental U. minor was occasionally planted as an ornamental urban tree. Augustine Henry wrote in 1913 that the U. minor planted in parks in Scotland were of French origin. More recently U. minor seed was imported to the UK from Italy. There are mature survivors in Edinburgh that are not the common U. minor cultivars (2015).
U. minor has been introduced to the southern hemisphere, notably Australasia and Argentina.
Notable trees
U. minor can live to a great age. An ancient field elm stood until recently in the village square of Metaxades, Thrace, Greece. Having abandoned their original village in 1286 after cholera outbreaks, the villagers re-founded it in the hills where a young elm was growing beside a spring. An elm (reputedly the original) and the fountain were the focal-point of the village until the late 20th century. The tallest recorded field elms in Greece were two specimens planted in 1650 beside the newly built church of the Archangels Michael and Gabriel, in Omali Voiou (Oμαλή Bοΐου) near Siatista, which, despite being open-grown trees, attained a height of 40 metres by the mid-20th century. The immemorial elm opposite the village square of Aidona in Thessaly, Greece which has been "listed" as a national "Monument of Nature", lost its crown in a recent storm (2009) and has now been pollarded; it is regenerating vigorously. A rare example of a centuries-old field elm that retains its heartwood and crown is the 360-year-old specimen in the village square of Strinylas, Corfu.
A tree said to be of similar age (200 cm d.b.h.) still stands (2013) in the city of Sliven, Bulgaria; other veterans are said to survive in the village of Samuilovo, 7 km from Sliven.
In France, a tree reputedly over 650 years old survived in the centre of Biscarrosse south of Bordeaux until the summer of 2010, when it finally succumbed to Dutch elm disease. Another veteran with a 6-metre girth survives at Bettange, France, close to the Belgian border, reputedly planted in 1593. Other wrecks include 'l'ormeau de Sully' in Villesèquelande near Carcassonne, "a magnificent tree supported by three metal props", said to have been planted in the early 17th C by the Duc de Sully,
A tree approximately 400 years old and 5.55 metres in girth grows in the town of Mergozzo in Piedmont, Italy. 'L'olmo di Mergozzo', like its French counterparts 'l'orme de Biscarosse' and 'l’orme de Bettange', is hollowed out by age, its life prolonged by lopping, while in Spain the elm in the Plaza del Olmo ("Elm square" in Spanish) in Navajas, Valencia, is 6.3 metres in girth; planted in 1636 it features on the town crest.
In England, large specimens once identified as U. minor subsp. minor, the narrow- or smooth-leafed elm, were once commonplace in the eastern counties before the advent of DED. The largest recorded tree in the UK grew at Great Amwell, Hertfordshire, measuring 40 m in height and 228 cm d.b.h. in 1911. Another famous specimen was the great elm that towered above its two siblings at the bottom of Long Melford Green, Long Melford, Suffolk, till the group succumbed to disease in 1978. The three "were survivors of a former clone of at least nine elms, one dating from 1757". The Long Melford elms were painted in 1940 by the watercolourist S. R. Badmin in his 'Long Melford Green on a Frosty Morning', now in the Victoria and Albert Museum. The largest known surviving trees in England are at East Coker, Somerset (30 m high, 95 cm d.b.h.), Termitts Farm near Hatfield Peverel, Essex (25 m high, 145 d.b.h.), and Melchbourne, Bedfordshire, (147 cm d.b.h.).
Subspecies, varieties, and former species sunk as U. minor
England
The name Ulmus minor subsp. minor was used by R. H. Richens for field elm that was not English elm, Cornish elm, Plot elm or Guernsey elm. Many publications, however, continue to use plain Ulmus minor for undifferentiated field elm; indeed Dr Max Coleman of Royal Botanic Garden Edinburgh argued in his 2002 paper 'British Elms' that there was no clear distinction between species and subspecies. Some authorities, among them Richens and Coleman, include English elm among varieties of field elm, Richens calling English elm U. minor var. vulgaris. Richens sank as undifferentiated U. minor certain local English forms such as U. minor 'Goodyeri', U. minor 'Hunnybunii', U. minor 'Sowerbyi', and U. minor 'Coritana'.
Eurasia
Henry's Ulmus nitens var. italica, 'Mediterranean Elm' (1913), distinguished by its 14 to 18 pairs of leaf-veins, was accepted, despite the wide source-area claimed for it ("Italy, Spain, Portugal and Algeria"), as U. carpinifolia var. italica Henry, by Krüssman (1984), who included a photograph of a specimen in Gisselfeld Park, Denmark. Bean (1988), however, considered it "a variety of rather dubious standing", and it was ignored by Richens (1983).
U. canescens Melville and U. boissieri Grudz. were both sunk as U. minor by Richens. The former is found throughout the eastern Mediterranean, including Palestine and Israel, and is distinguished by its leaves, densely downy on the underside when mature. The latter is a little-known tree found in Iran, in the Zagros forests and the Kerman / Kermanshah area.
Green and Richens also sank U. minor var. suberosa (Moench) Rehder - the so-called 'Cork-barked elm', korkulme (Germany) or wiąz korkowa (Poland), as a genetically random, maritime or juvenile form of U. minor, insufficiently differentiated to merit varietal status, its name a relic of taxonomic conservatism.
Cultivars
Numerous cultivars have been raised in Europe since the 18th century, although many are now probably either extinct owing to the ravages of Dutch elm disease, or survive unrecognized in sucker form:
Hybrids
The tree's natural range generously overlaps that of wych elm Ulmus glabra to the north, and readily hybridizes with it to produce the so-called 'Dutch elm' Ulmus × hollandica.
In Spain and Italy Ulmus minor has naturally hybridized with Siberian elm U. pumila. In Spain U. pumila was introduced in the 16th century and has since spread widely, contributing to conservation concerns for U. minor. In Italy U. pumila was introduced in the 1930s; research is ongoing into the extent of its hybridisation with U. minor. The resulting hybrid has not yet been given a formal botanical name, though there are cultivated forms such as 'Recerta' and 'Fiorente' (see 'Hybrid cultivars').
Ulmus × hollandica
Ulmus davidiana var. japonica × U. minor
U. minor × U. pumila
Hybrid cultivars
U. minor hybridises naturally with U. glabra, producing elms of the Ulmus × hollandica group, from which there have arisen a number of cultivars:
The tree has featured strongly in artificial hybridization experiments in Europe and to a lesser extent in the United States. The hybrid Ulmus davidiana var. japonica × U. minor was raised at the Arnold Arboretum before 1924. Most of the European research was based at Wageningen in the Netherlands until 1992, whence a number of hybrid cultivars have been commercially released since 1960. The earlier trees were raised in response to the initial Dutch elm disease pandemic that afflicted Europe after the First World War, and were to prove vulnerable to the much more virulent strain of the disease that arrived in the late 1960s. However, further research eventually produced several trees highly resistant to disease which were released after 1989.
Arno, Clusius, Columella, Commelin, Den Haag, Fiorente, Frontier, Fuente Umbria, Groeneveld, Homestead, Lobel, Nanguen = , Pioneer, Plantyn, Plinio, Recerta, San Zanobi, Toledo, Urban, Wanoux =
In art
The elms by Willy Lott's Cottage and Flatford Mill, Suffolk, in Constable's paintings and drawings were, according to Richens, "smooth-leaved elm" (U. minor), though the hedgerow elms in his Dedham Vale and East Bergholt landscape-paintings and drawings were otherwise "most probably East Anglian hybrid elms ... such as still grow in the same hedges".
Accessions
North America
United States National Arboretum, Washington, D.C., US. Acc. nos. 12852, 64382.
Europe
Arboretum de La Petite Loiterie , Monthodon, France. No details available
Cambridge Botanic Garden , University of Cambridge, UK. No accession details available.
Dubrava Arboretum, Lithuania. No details available.
Grange Farm Arboretum, Sutton St James, Spalding, Lincolnshire, UK. Acc. no. not known.
Linnaean Gardens of Uppsala, Finland. Acc. no. 1930-1013.
Royal Botanic Garden Edinburgh, UK. Acc. nos. 19699368, 16899359, 19699365.
Royal Botanic Gardens, Kew, UK. Acc. no. not known.
Sir Harold Hillier Gardens, Hampshire, UK. Acc. no. 2001-0188, 3 specimens collected in Iran, 2000.
Strona Arboretum, University of Life Sciences, Warsaw, Poland. No details available.
Australasia
Eastwoodhill Arboretum , Gisborne, New Zealand. 2 trees, details not known.
Nurseries
North America
None known
Europe
Eggleston Hall Gardens , Eggleston, Barnard Castle, County Durham, UK
Firecrest Tree & Shrub Nursery , Woodbridge, Suffolk, UK
Lorenz von Ehren , Hamburg, Germany
Trees & Hedges , Heathfield, East Sussex, UK
UmbraFlor , Spello, Italy
| Biology and health sciences | Rosales | Plants |
4814988 | https://en.wikipedia.org/wiki/Najash | Najash | Najash is an extinct genus of basal snake from the Late Cretaceous Candeleros Formation of Patagonia. Like a number of other Cretaceous and living snakes it retained hindlimbs, but Najash is unusual in having well-developed legs that extend outside the rib cage, and a pelvis connected to the spine.
Discovery and description
Fossils of Najash were found in the terrestrial Candeleros Formation, in Rio Negro Province, Argentina, and date to roughly 90 million years ago. The skull and spine of Najash show primitive features that resemble other Cretaceous snakes, such as Dinilysia patagonica and Madtsoiidae. Also, several characteristics of the neck and tail of Najash and Dinilysia patagonica show how the body plan of snakes evolved from a lizard-like ancestor.
Najash had not lost its sacrum, the pelvic bone composed of several fused vertebrae, nor its pelvic girdle, which are absent in modern snakes, and in all other known fossil snakes as well. Nearly all phylogenetic analyses place Najash as an early offshoot of the snake tree, outside of all living snakes.
| Biology and health sciences | Prehistoric squamates | Animals |
1263998 | https://en.wikipedia.org/wiki/Slow%20loris | Slow loris | Slow lorises are a group of several species of nocturnal strepsirrhine primates that make up the genus Nycticebus. Found in Southeast Asia and nearby areas, they range from Bangladesh and Northeast India in the west to the Sulu Archipelago in the Philippines in the east, and from Yunnan province in China in the north to the island of Java in the south.
Although many previous classifications recognized as few as a single all-inclusive species, there are now at least eight that are considered valid: the Sunda slow loris (N. coucang), Bengal slow loris (N. bengalensis), Javan slow loris (N. javanicus), Philippine slow loris (N. menagensis), Bangka slow loris (N. bancanus), Bornean slow loris (N. borneanus), Kayan River slow loris (N. kayan) and Sumatran slow loris (N. hilleri). A ninth species, the pygmy slow loris (X. pygmaeus), was recently moved to the new genus Xanthonycticebus. After the pygmy slow loris, the group's closest relatives are the slender lorises of southern India and Sri Lanka. Their next closest relatives are the African lorisids, the pottos, false pottos, and angwantibos. They are less closely related to the remaining lorisoids (the various types of galago), and more distantly to the lemurs of Madagascar. Their evolutionary history is uncertain since their fossil record is patchy and molecular clock studies have given inconsistent results.
Slow lorises have a round head, a narrow snout, large eyes, and a variety of distinctive coloration patterns that are species-dependent. Their arms and legs are nearly equal in length, and their torso is long and flexible, allowing them to twist and extend to nearby branches. The hands and feet of slow lorises have several adaptations that give them a pincer-like grip and enable them to grasp branches for long periods of time. Slow lorises have a toxic bite, a trait rare among mammals and unique among the primates. The toxin is obtained by licking a sweat gland on their arm, and the secretion is activated by mixing with saliva. Their toxic bite, once thought to be primarily a deterrent to predators, has been discovered to be primarily used in disputes within the species.
The secretion from the arm contains a chemical related to cat allergen, but may be augmented by secondary toxins from the diet in wild individuals. Slow lorises move slowly and deliberately, making little or no noise, and when threatened, they stop moving and remain motionless. Their only documented predators—apart from humans—include snakes, changeable hawk-eagles and orangutans, although cats, viverrids and sun bears are suspected. Little is known about their social structure, but they are known to communicate by scent marking. Males are highly territorial. Slow lorises reproduce slowly, and the infants are initially parked on branches or carried by either parent. They are omnivores, eating small animals, fruit, tree gum, and other vegetation.
Each of the slow loris species that had been identified prior to 2012 is listed as either "Vulnerable" or "Endangered" on the IUCN Red List. The three newest species are yet to be evaluated, but they arise from (and further reduce the ranks of) what was thought to be a single "vulnerable" species. All four of these are expected to be listed with at least the same, if not a higher-risk, conservation status. All slow lorises are threatened by the wildlife trade and habitat loss. Their habitat is rapidly disappearing and becoming fragmented, making it nearly impossible for slow lorises to disperse between forest fragments; unsustainable demand from the exotic pet trade and from traditional medicine has been the greatest cause for their decline.
Taxonomy and systematics
Although many previous classifications recognized as few as a single all-inclusive species, there are now at least eight that are considered valid:
Other than the pygmy slow loris in sister genus Xanthonycticebus, the group's closest relatives are the slender lorises of southern India and Sri Lanka. Their next closest relatives are the African lorisids, the pottos, false pottos, and angwantibos. They are less closely related to the remaining lorisoids (the various types of galago), and more distantly to the lemurs of Madagascar. Their evolutionary history is uncertain since their fossil record is patchy and molecular clock studies have given inconsistent results.
Evolutionary history
Slow lorises (genus Nycticebus) are strepsirrhine primates and are related to other living lorisoids, such as the pygmy slow loris (Xanthonycticebus), slender lorises (Loris), pottos (Perodicticus), false pottos (Pseudopotto), angwantibos (Arctocebus), and galagos (family Galagidae), and to the lemurs of Madagascar. They are most closely related to the pygmy slow loris, followed by the slender lorises of South Asia, the angwantibos, pottos and false pottos of Central and West Africa. Lorisoids are thought to have evolved in Africa, where most living species occur; later, one group may have migrated to Asia and evolved into the slender and slow lorises of today.
Lorises first appear in the Asian fossil record in the Miocene, with records in Thailand around 18 million years ago (mya) and in Pakistan 16 mya. The Thai record is based on a single tooth that most closely resembles living slow lorises and that is tentatively classified as a species of Nycticebus. The species is named ? Nycticebus linglom, using open nomenclature (the preceding "?" indicates the tentative nature of the assignment).
Several lorises are found in the Siwalik deposits of Pakistan, dating to 16 to 8 mya, including Nycticeboides and Microloris. Most are small, but an unnamed form dating to 15–16 mya is comparable in size to the largest living slow lorises. Molecular clock analysis suggests that slow lorises may have started evolving into distinct species about 10 mya. They are thought to have reached the islands of Sundaland when the Sunda Shelf was exposed at times of low sea level, creating a land bridge between the mainland and islands off the coast of Southeast Asia.
Discovery and taxonomy
The earliest known mention of a slow loris in scientific literature is from 1770, when Dutchman Arnout Vosmaer (1720–1799) described a specimen of what we know today as N. bengalensis that he had received two years earlier. The French naturalist Georges-Louis Leclerc, Comte de Buffon, later questioned Vosmaer's decision to affiliate the animal with sloths, arguing that it was more closely aligned with the lorises of Ceylon (now Sri Lanka) and Bengal. The word "loris" was first used in 1765 by Buffon as a close equivalent to a Dutch name, loeris. This etymology was later supported by the physician William Baird in the 1820s, who noted that the Dutch word loeris signified "a clown".
In 1785, the Dutch physician and naturalist Pieter Boddaert was the first to officially describe a species of slow loris using the name Tardigradus coucang. This species was based on the "tailless maucauco" described by Thomas Pennant in 1781, which is thought to have been based on a Sunda slow loris, and on Vosmaer's description of a Bengal slow loris. Consequently, there has been some disagreement over the identity of Tardigradus coucang; currently the name is given to the Sunda slow loris. The next slow loris species to be described was Lori bengalensis (currently Nycticebus bengalensis), named by Bernard Germain de Lacépède in 1800.
In 1812, Étienne Geoffroy Saint-Hilaire named the genus Nycticebus, naming it for its nocturnal behavior. The name derives from the , genitive form of (, "night"), and (, "monkey"). Geoffroy also named Nycticebus javanicus in this work. Later 19th-century authors also called the slow lorises Nycticebus, but most used the species name tardigradus (given by Linnaeus in 1758 in the 10th edition of Systema Naturæ) for slow lorises, until mammalogists Witmer Stone and James A. G. Rehn clarified in 1902 that Linnaeus's name actually referred to a slender loris.
Several more species were named around 1900, including Nycticebus menagensis (originally Lemur menagensis) by Richard Lydekker in 1893 and Nycticebus pygmaeus by John James Lewis Bonhote in 1907. However, in 1939 Reginald Innes Pocock consolidated all slow lorises into a single species, N. coucang, and in his influential 1953 book Primates: Comparative Anatomy and Taxonomy, primatologist William Charles Osman Hill also followed this course. In 1971 Colin Groves recognized the pygmy slow loris (N. pygmaeus) as a separate species, and divided N. coucang into four subspecies, while in 2001 Groves opined there were three species (N. coucang, N. pygmaeus, and N. bengalensis), and that N. coucang had three subspecies (Nycticebus coucang coucang, N. c. menagensis, and N. c. javanicus).
In 2006, the Bornean slow loris was elevated to the species level (as Nycticebus menagensis) based on molecular analysis of DNA sequences of the D-loop and the cytochrome b gene. In 2008, Groves and Ibnu Maryanto confirmed the promotion of the fifth species, the Javan slow loris, to species status, a move that had been suggested in previous studies from 2000. They based their decision on an analysis of cranial morphology and characteristics of pelage. Species differentiation was based largely on differences in morphology, such as size, fur color, and head markings.
To help clarify species and subspecies boundaries, and to establish whether morphology-based classifications were consistent with evolutionary relationships, the phylogenetic relationships within the genus Nycticebus were investigated by Chen and colleagues using DNA sequences derived from the mitochondrial markers D-loop and cytochrome b. Previous molecular analyses using karyotypes, restriction enzymes, and DNA sequences were focused on understanding the relationships between a few species, not the phylogeny of the entire genus. The analyses published in 2006 by Chen and colleagues' proved inconclusive, although one test suggested that N. coucang and N. bengalensis apparently share a closer evolutionary relationship with each other than with members of their own species, possibly due to introgressive hybridization since the tested individuals of these two taxa originated from a region of sympatry in southern Thailand. This hypothesis was corroborated by a 2007 study that compared the variations in mitochondrial DNA sequences between N. bengalensis and N. coucang, and suggested that there has been gene flow between the two species.
In 2012, two taxonomic synonyms (formerly recognized as subspecies) of N. menagensis—N. bancanus and N. borneanus—were elevated to species status, and a new species—N. kayan—was also distinguished from the same. Rachel Munds, Anna Nekaris and Susan Ford based these taxonomic revisions on distinguishable facial markings. With that, the N. menagensis species complex that had been collectively known as the Bornean slow loris became four species: the Philippine slow loris (N. menagensis), the Bornean slow loris (N. borneanus), the Bangka slow loris (N. bancanus), and the Kayan River slow loris (N. kayan).
Nekaris and Nijman (2022) combined morphological, behavioural, karyotypical and genetic data and suggested that the pygmy slow loris is best placed in its own genus, Xanthonycticebus.
Anatomy and physiology
Slow lorises have a round head because their skull is shorter than in other living strepsirrhine. Like other lorisids, their snout does not taper towards the front of the face as it does in lemurs, making the face appear less long and pointed. Compared with the slender lorises, the snout of the slow loris is even less pointed. As with other members of Lorisidae, its interorbital distance is shorter than in lemurs. The skull has prominent crests (ridges of bone). A distinguishing feature of the slow loris skull is that the occipital bone is flattened and faces backward. The foramen magnum (hole through which the spinal cord enters) faces directly backward. The brains of slow lorises have more folds (convolutions) than the brains of galagos.
The ears are small, sparsely covered in hair, and hidden in the fur. Similar to the slender lorises, the fur around and directly above the eyes is dark. Unlike the slender lorises, however, the white stripe that separates the eye rings broadens both on the tip of the nose and on the forehead while also fading out on the forehead. Like other strepsirrhine primates, the nose and lip are covered by a moist skin called the rhinarium ("wet nose"), which is a sense organ.
The eyes of slow lorises are forward-facing, which gives stereo vision. Their eyes are large and possess a reflective layer, called the tapetum lucidum, that improves low-light vision. It is possible that this layer blurs the images they see, as the reflected light may interfere with the incoming light. Slow lorises have monochromatic vision, meaning they see in shades of only one color. They lack the opsin gene that would allow them to detect short wavelength light, which includes the colors blue and green.
The dental formula of slow lorises is , meaning that on each side of the mouth there are two upper (maxillary) and lower (mandibular) incisors, one upper and lower canine tooth, three upper and lower premolars, and three upper and lower molars, giving a total of 36 permanent teeth. As in all other crown strepsirrhines, their lower incisors and canine are procumbent (lie down and face outwards), forming a toothcomb, which is used for personal and social grooming and feeding. The toothcomb is kept clean by the sublingua or "under-tongue", a specialized structure that acts like a toothbrush to remove hair and other debris. The sublingua extends below the tip of the tongue and is tipped with keratinized, serrated points that rake between the front teeth.
Slow lorises have relatively large maxillary canine teeth, their inner (mesial) maxillary incisors are larger than the outer (distal) maxillary incisors, and they have a diastema (gap) between the canine and the first premolar. The first mandibular premolar is elongated, and the last molar has three cusps on the crown, the shortest of which is near the back. The bony palate (roof of the mouth) only goes as far back as the second molar.
Slow lorises range in weight from the Bornean slow loris at to as much as for the Bengal slow loris. Slow lorises have stout bodies, and their tails are only stubs and hidden beneath the dense fur. Their combined head and body lengths vary by species, but range from between all species. The trunk is longer than in other living strepsirrhines because they have 15–16 thoracic vertebrae, compared to 12–14 in other living strepsirrhines. This gives them greater mobility when twisting and extending towards nearby branches. Their other vertebrae include seven cervical vertebrae, six or seven lumbar vertebrae, six or seven sacral vertebrae, and seven to eleven caudal vertebrae.
Unlike galagos, which have longer legs than arms, slow lorises have arms and legs of nearly equal length. Their intermembral index (ratio of arm to leg length) averages 89, indicating that their forelimbs are slightly shorter than their hind limbs. As with the slender lorises, their arms are slightly longer than their body, but the extremities of slow lorises are more stout.
Slow lorises have a powerful grasp with both their hands and feet due to several specializations. They can tightly grasp branches with little effort because of a special muscular arrangement in their hands and feet, where the thumb diverges at nearly 180° from the rest of the fingers, while the hallux (big toe) ranges between being perpendicular and pointing slightly backwards. The toes have a large flexor muscle that originates on the lower end of the thigh bone, which helps to impart a strong grasping ability to the hind limbs.
The second digit of the hand is short compared to the other digits, while on the foot, the fourth toe is the longest. The sturdy thumb helps to act like a clamp when digits three, four, and five grasp the opposite side of a tree branch. This gives their hands and feet a pincer-like appearance. The strong grip can be held for hours without losing sensation due to the presence of a rete mirabile (network of capillaries), a trait shared among all lorises. Both slender and slow lorises have relatively short feet. Like nearly all lemuriforms, they have a grooming claw on the second toe of each foot.
Slow lorises have an unusually low basal metabolic rate, about 40% of the typical value for placental mammals of their size, comparable to that of sloths. Since they consume a relatively high-calorie diet that is available year-round, it has been proposed that this slow metabolism is due primarily to the need to eliminate toxic compounds from their food. For example, slow lorises can feed on Gluta bark, which may be fatal to humans.
Distribution and diversity
Slow lorises are found in South and Southeast Asia. Their collective range stretches from Northeast India through Indochina, east to the Sulu Archipelago (the small, southern islands of the Philippines), and south to the island of Java (including Borneo, Sumatra, and many small nearby islands). They are found in India (Northeastern states), China (Yunnan province), Laos, Vietnam, Cambodia, Bangladesh, Burma, Thailand, Malaysia, the Philippines, Indonesia, Brunei, and Singapore.
There are currently seven recognized species. The Bornean slow loris (N. menagensis), found on Borneo and nearby islands, including the Sulu Archipelago, and in 2012 was split into four distinct species (adding N. bancanus, N. borneanus, and N. kayan). The Javan slow loris (N. javanicus) is only found on the island of Java in Indonesia. The Sunda slow loris (N. coucang) occurs on Sumatra and the Malay Peninsula, including Singapore and southern Thailand (the Isthmus of Kra). The Bengal slow loris (N. bengalensis) has the largest distribution of all the slow lorises and can be found in Bangladesh, Cambodia, southern China, Northeast India, Laos, Burma, Thailand, and Vietnam.
Slow lorises range across tropical and subtropical regions and are found in primary and secondary rainforests, as well as bamboo groves and mangrove forests. They prefer forests with high, dense canopies, although some species have also been found in disturbed habitats, such as cacao plantations and mixed-crop home gardens. Due largely to their nocturnal behavior and the subsequent difficulties in accurately quantifying abundance, data about the population size or distribution patterns of slow lorises is limited. In general, encounter rates are low; a combined analysis of several field studies involving transect surveys conducted in South and Southeast Asia determined encounter rates ranging from as high as 0.74 lorises per kilometer for N. coucang to as low as 0.1 lorises per kilometer for N. bengalensis.
Behavior and ecology
Little is known about the social structure of slow lorises, but they generally spend most of the night foraging alone. Individuals sleep during the day, usually alone but occasionally with other slow lorises. Home ranges of adults may significantly overlap, and those of males are generally larger than those of females. In the absence of direct studies of the genus, primatologist Simon Bearder speculated that slow loris social behavior is similar to that of the potto, another nocturnal primate.
Such a social system is distinguished by a lack of matriarchy and by factors that allow the slow loris to remain inconspicuous and minimize energy expenditure. Vocal exchanges and alarm calls are limited; scent marking with urine is the dominant form of communication. Adult males are highly territorial and are aggressive towards other males. Vocalizations include an affiliative (friendly) call krik, and a louder call resembling a crow's caw. When disturbed, slow lorises can also produce a low buzzing hiss or growl. To make contact with other individuals, they emit a single high-pitched rising tone, and females use a high whistle when in estrus.
Slow lorises are slow and deliberate climbers, and often hold on to branches with three of their four limbs. To move between trees, they carefully grip the terminal branches of the neighboring tree and pull themselves across the small gap. They will also grip branches with only their hind feet, lift themselves upright, and quickly launch forward with their hands to catch prey.
Due to their slow movement, all lorises, including the slow lorises, have a specially adapted mechanism for defense against predation. Their slow, deliberate movement hardly disturbs the vegetation and is almost completely silent. Once disturbed, they immediately stop moving and remain motionless. In Indonesia, slow lorises are called malu malu or "shy one" because they freeze and cover their face when spotted.
If cornered, they may adopt a defensive posture by curling up and lunging at the predator. The Acehnese name, buah angin ("wind monkey"), refers to their ability to "fleetingly but silently escape". Little is known about the predation of slow lorises. Documented predators include snakes, the changeable hawk-eagle (Nisaetus cirrhatus), and Sumatran orangutans (Pongo abelii). Other potential predators include cats, sun bears (Helarctos malayanus), binturongs (Arctictis binturong), and Asian palm civets.
Slow lorises produce a secretion from their brachial gland (a scent gland on the upper arm near the axilla) that is licked and mixed with their saliva. In tests, three predators—binturongs, clouded leopards (Neofelis nebulosa), and sun bears—retreated or showed other signs of displeasure when presented with cotton swabs anointed with a mixture of the toxic secretion and the saliva, whereas the toxic secretion alone generated mild interest. Before stashing their offspring in a secure location, female slow lorises will lick their brachial glands, and then groom their young with their toothcomb, depositing the toxin on their fur. When threatened, slow lorises may also lick their brachial glands and bite their aggressors, delivering the toxin into the wounds. Slow lorises can be reluctant to release their bite, which is likely to maximize the transfer of toxins. This toxic bite is a rare trait among mammals and unique to lorisid primates. It may also be used for defense against other slow lorises and parasites. According to Nekaris, this adaptation—along with vocalizations, movement, and coloration patterns similar to those of true cobras—may have evolved through Müllerian mimicry to protect slow lorises when they need to move across the ground due to breaks in the canopy.
The secretion from the brachial gland of captive slow lorises is similar to the allergen in cat dander, hence the secretions may merely elicit an allergic reaction, not toxicosis. Loris bites cause a painful swelling, and the single case of human death reported in the scientific literature was believed to have resulted from anaphylactic shock.
To protect itself, the Slow loris has also been observed to rub the venom on its fur to chemically defend itself from predators.
Studies suggest that slow lorises are polygynandrous. Infants are either parked on branches while their parents find food or else are carried by one of the parents. Due to their long gestations (about six months), small litter sizes, low birth weights, long weaning times (three to six months), and long gaps between births, slow loris populations have one of the slowest growth rates among mammals of similar size. Pygmy slow lorises are likely to give birth to twins—from 50% to 100% of births, depending on the study; in contrast, this phenomenon is rare (3% occurrence) in Bengal slow lorises. A seven-year study of captive-bred pygmy slow lorises showed a skewed sex distribution, with 1.68 males born for every 1 female.
Breeding may be continuous throughout the year. Copulation often occurs while suspended with the hands and feet clinging to horizontal branches for support. In captive Sunda slow lorises, mating primarily occurs between June and mid-September, with the estrus cycle lasting 29 to 45 days and estrus lasting one to five days. Likewise, gestation lasts 185 to 197 days, and the young weigh between at birth. Females reach sexual maturity at 18 to 24 months, while males are capable of reproducing at 17 months. However, the fathers become hostile towards their male offspring after 12 to 14 months and will chase them away. In captivity, they can live 20 or more years.
Diet
Slow lorises are omnivores, eating insects and other arthropods, small birds and reptiles, eggs, fruits, gums, nectar and miscellaneous vegetation. A 1984 study of the Sunda slow loris indicated that its diet consists of 71% fruit and gums, and 29% insects and other animal prey. A more detailed study of another Sunda slow loris population in 2002 and 2003 showed different dietary proportions, consisting of 43.3% gum, 31.7% nectar, 22.5% fruit, and just 2.5% arthropods and other animal prey. The most common dietary item was nectar from flowers of the Bertram palm (Eugeissona tristis). The Sunda slow loris eats insects that other predators avoid due to their repugnant taste or smell.
Preliminary results of studies on the pygmy slow loris indicate that its diet consists primarily of gums and nectar (especially nectar from Saraca dives flowers), and that animal prey makes up 30–40% of its diet. However, one 2002 analysis of pygmy slow loris feces indicated that it contained 98% insect remains and just 2% plant remains. The pygmy slow loris often returns to the same gum feeding sites and leaves conspicuous gouges on tree trunks when inducing the flow of exudates. Slow lorises have been reported gouging for exudates at heights ranging from to as much as ; the gouging process, whereby the loris repetitively bangs its toothcomb into the hard bark, may be loud enough to be heard up to away. The marks remaining after gouging can be used by field workers to assess loris presence in an area.
Captive pygmy slow lorises also make characteristic gouge marks in wooden substrates, such as branches. It is not known how the sympatric pygmy and Bengal slow lorises partition their feeding niches. The plant gums, obtained typically from species in the family Fabaceae (peas), are high in carbohydrates and lipids, and can serve as a year-around source of food, or an emergency reserve when other preferred food items are scarce. Several anatomical adaptations present in slow lorises may enhance their ability to feed on exudates: a long narrow tongue to make it easier to reach gum stashed in cracks and crevices, a large cecum to help the animal digest complex carbohydrates, and a short duodenum to help quickly pass potentially toxic exudates. Slow lorises can use both hands to eat while hanging upside down from a branch. They spend about 20% of their nightly activities feeding.
In culture
Beliefs about slow lorises and their use in traditional practices are deep-rooted and go back at least 300 years, if not earlier based on oral traditions. In the late 19th and early 20th centuries, it was reported that the people from the interior of Borneo believed that slow lorises were the gatekeepers for the heavens and that each person had a personal slow loris waiting for them in the afterlife. More often, however, slow lorises are used in traditional medicine or to ward off evil. The following passage from an early textbook about primates is indicative of the superstitions associated with slow lorises: Many strange powers are attributed to this animal by the natives of the countries it inhabits; there is hardly an event in life to man, woman or child, or even domestic animals, that may not be influenced for better or worse by the Slow Loris, alive or dead, or by any separate part of it, and apparently one cannot usually tell at the time, that one is under supernatural power. Thus a Malay may commit a crime he did not premeditate, and then find that an enemy had buried a particular part of a Loris under his threshold, which had, unknown to him, compelled him to act to his own disadvantage. ... [a slow loris's] life is not a happy one, for it is continually seeing ghosts; that is why it hides its face in its hands.
In the Mondulkiri Province of Cambodia, hunters believe that lorises can heal their own broken bones immediately after falling from a branch so that they can climb back up the tree. They also believe that slow lorises have medicinal powers because they require more than one hit with a stick to die. In the province of North Sumatra, the slow loris is thought to bring good luck if it is buried under a house or a road. In the same province, slow loris body parts were used to place curses on enemies. In Java, it was thought that putting a piece of its skull in a water jug would make a husband more docile and submissive, just like a slow loris in the daytime. More recently, researchers have documented the belief that the consumption of loris meat was an aphrodisiac that improves "male power". The gall bladder of the Bengal slow loris has historically been used to make ink for tattoos by the village elders in Pursat and Koh Kong Provinces of Cambodia. Loris wine is a traditional Cambodian medicine supposed to alleviate the pain of childbirth, made from a mixture of loris bodies and rice wine.
Conservation
The two greatest threats to slow lorises are deforestation and the wildlife trade. Slow lorises have lost a significant amount of habitat, with habitat fragmentation isolating small populations and obstructing biological dispersal. However, despite the lost habitat, their decline is most closely associated with unsustainable trade, either as exotic pets or for traditional medicine.
Each of the slow loris species that had been identified prior to 2012 are currently listed as either "Vulnerable" or "Endangered" by the International Union for Conservation of Nature (IUCN) on their Red List. When they were all considered a single species, imprecise population data together with their regular occurrence in Southeast Asian animal markets combined to incorrectly suggest that slow lorises were common. This manifested as incorrect Red List assessments of "Least Concern" as recently as 2000. The three newest species are yet to be evaluated by the IUCN, although each was once thought to be subpopulations of the Bornean slow loris—which was evaluated as "Vulnerable" in 2008. With this division of its range and population, the Bornean slow loris and the three new species face a higher risk of extinction than before.
Since 2007, all slow loris species have been protected from commercial international trade under Appendix I of CITES. Furthermore, local trade is illegal because every nation in which they occur naturally has laws protecting them. Despite their CITES Appendix I status and local legal protection, slow lorises are still threatened by both local and international trade due to problems with enforcement. Surveys are needed to determine existing population densities and habitat viability for all species of slow loris. Connectivity between protected areas is important for slow lorises because they are not adapted to dispersing across the ground over large distances.
Populations of Bengal and Sunda slow lorises are not faring well in zoos. Of the 29 captive specimens in North American zoos in 2008, several are hybrids that cannot breed, while most are past their reproductive years. The last captive birth for these species in North America was in 2001 in San Diego. Pygmy slow lorises are doing better in North American zoos; from the late 1980s (when they were imported) to 2008, the population grew to 74 animals, with most of them born at the San Diego Zoo.
Wildlife trade
Until the 1960s, the hunting of slow lorises was sustainable, but due to growing demand, decreased supply, and the subsequent increased value of the marketed wildlife, slow lorises have been overexploited and are in decline. With the use of modern technology, such as battery-powered searchlights, slow lorises have become easier to hunt because of their eyeshine. Traditional medicine made from loris parts is thought to cure many diseases, and the demand for this medicine from wealthy urban areas has replaced the subsistence hunting traditionally performed in poor rural areas. A survey by primatologist Anna Nekaris and colleagues (2010) showed that these belief systems were so strong that the majority of respondents expressed reluctance to consider alternatives to loris-based medicines.
Slow lorises are sold locally at street markets but are also sold internationally over the Internet and in pet stores. They are especially popular or trendy in Japan, particularly among women. The reasons for their popularity, according to the Japan Wildlife Conservation Society, are that "they're easy to keep, they don't cry, they're small, and just very cute."
Common misconceptions
Because of their "cuteness", videos of pet slow lorises are some of the most frequently watched animal-related viral videos on YouTube. By March 2011, a newly posted video of a slow loris holding a cocktail umbrella had been viewed more than two million times, while an older video of a slow loris being tickled had been viewed more than six million times. According to Nekaris, these videos are misunderstood by most people who watch them, since most do not realize that it is illegal in most countries to own them as pets and that the slow lorises in the videos are only docile because that is their passive defensive reaction to threatening situations.
Despite frequent advertisements by pet shops in Japan, the World Conservation Monitoring Centre reported that only a few dozen slow lorises were legally imported in 2006, suggesting frequent smuggling. Slow lorises are also smuggled to China, Taiwan, Europe, Russia, the United States, and Saudi Arabia for use as pets.
Within their countries of origin, slow lorises are very popular pets, particularly in Indonesia. They are seen as a "living toy" for children by local people or are bought out of pity by Western tourists or expatriates. Neither local nor foreign buyers usually know anything about these primates, their endangered status, or that the trade is illegal. According to National Geographic, slow lorises are protected by both local laws in southern Asia and by the Convention on International Trade in Endangered Species (CITES). Furthermore, few know about their strong odor or their painful bite, which may lead to anaphylaxis in some cases. According to data compiled from monthly surveys and interviews with local traders, nearly a thousand locally-sourced slow lorises exchanged hands in the Medan bird market in North Sumatra during the late first decade of the 21st century.
International trade usually involves a high mortality rate during transit, between 30% and 90%. Slow lorises also experience many health problems due to both local and international trade. In order to give the impression that the primates are tame and appropriate pets for children, to protect people from their potentially toxic bite, or to deceive buyers into thinking the animal is a baby, animal dealers either pull the front teeth with pliers or wire cutters or cut them off with nail clippers. This results in severe bleeding, which sometimes causes shock or death.
Dental infection is common and is fatal in 90% of cases. Without their teeth, the animals can no longer fend for themselves in the wild and must remain in captivity for life. The slow lorises found in animal markets are usually underweight and malnourished and have had their fur dyed, which complicates species identification at rescue centers. As many as 95% of the slow lorises rescued from the markets die of dental infection or improper care.
As part of the trade, infants are pulled prematurely from their parents, leaving them unable to remove their own urine, feces, and oily skin secretions from their fur. Slow lorises have a special network of blood vessels in their hands and feet, which makes them vulnerable to cuts when pulled from the wire cages they are kept in. Slow lorises are also stress-sensitive and do not thrive in captivity. Common health problems seen in pet slow lorises include undernourishment, tooth decay, diabetes, obesity, and kidney failure. Infection, stress, pneumonia, and poor nutrition lead to high death rates among pet lorises. Pet owners also fail to provide proper care because they are usually asleep when the nocturnal pet is awake.
| Biology and health sciences | Primates | null |
1264072 | https://en.wikipedia.org/wiki/Capella | Capella | Capella is the brightest star in the northern constellation of Auriga. It has the Bayer designation α Aurigae, which is Latinised to Alpha Aurigae and abbreviated Alpha Aur or α Aur. Capella is the sixth-brightest star in the night sky, and the third-brightest in the northern celestial hemisphere after Arcturus and Vega. A prominent object in the northern winter sky, it is circumpolar to observers north of 44°N. Its name meaning "little goat" in Latin, Capella depicted the goat Amalthea that suckled Zeus in classical mythology. Capella is relatively close, at from the Sun. It is one of the brightest X-ray sources in the sky, thought to come primarily from the corona of Capella Aa.
Although it appears to be a single star to the naked eye, Capella is actually a quadruple star system organized in two binary pairs, made up of the stars Capella Aa, Capella Ab, Capella H and Capella L. The primary pair, Capella Aa and Capella Ab, are two bright-yellow giant stars, both of which are around 2.5 times as massive as the Sun. The secondary pair, Capella H and Capella L, are around 10,000 astronomical units (AU) from the first and are two faint, small and relatively cool red dwarfs.
Capella Aa and Capella Ab have exhausted their core hydrogen, and cooled and expanded, moving off the main sequence. They are in a very tight circular orbit about 0.74 AU apart, and orbit each other every 104 days. Capella Aa is the cooler and more luminous of the two with spectral class K0III; it is 78.7 ± 4.2 times the Sun's luminosity and 11.98 ± 0.57 times its radius. An aging red clump star, it is fusing helium to carbon and oxygen in its core. Capella Ab is slightly smaller and hotter and of spectral class G1III; it is 72.7 ± 3.6 times as luminous as the Sun and 8.83 ± 0.33 times its radius. It is in the Hertzsprung gap, corresponding to a brief subgiant evolutionary phase as it expands and cools to become a red giant. Several other stars in the same visual field have been catalogued as companions but are physically unrelated.
Nomenclature
α Aurigae (Latinised to Alpha Aurigae) is the star system's Bayer designation. It also has the Flamsteed designation 13 Aurigae. It is listed in several multiple star catalogues as ADS 3841, CCDM J05168+4559, and WDS J05167+4600. As a relatively nearby star system, Capella is listed in the Gliese-Jahreiss Catalogue with designations GJ 194 for the bright pair of giants and GJ 195 for the faint pair of red dwarfs.
The traditional name Capella is Latin for (small) female goat; the alternative name Capra was more commonly used in classical times. It is the translation of the Greek star name Aἴξ (aix) meaning "the Goat". As the sound of the Greek term for the goat (aἴξ) is similar to the sound of the name for the Aegaean Sea, this star has been used for weather rules and determining the seasonal wind direction. In 2016, the International Astronomical Union organized a Working Group on Star Names (WGSN) to catalogue and standardize proper names for stars. The WGSN's first bulletin of July 2016 included a table of the first two batches of names approved by the WGSN; which included Capella for this star. It is now so entered in the IAU Catalog of Star Names. The catalogue of star names lists Capella as applying to the star α Aurigae Aa.
Observational history
Capella was the brightest star in the night sky from 210,000 years ago to 160,000 years ago, at about −1.8 in apparent magnitude. At −1.1, Aldebaran was brightest before this period; it and Capella were situated rather close to each other in the sky and approximated boreal pole stars at the time.
Capella is thought to be mentioned in an Akkadian inscription dating to the 20th century BC. Its goat-associated symbolism dates back to Mesopotamia as a constellation called "GAM", "Gamlum" or "MUL.GAM" in the 7th-century BC document MUL.APIN. GAM represented a scimitar or crook and may have represented the star alone or the constellation of Auriga as a whole. Later, Bedouin astronomers created constellations that were groups of animals, where each star represented one animal. The stars of Auriga comprised a herd of goats, an association also present in Greek mythology. It is sometimes called the Shepherd's Star in English literature. Capella was seen as a portent of rain in classical times.
Building J of the pre-Columbian site Monte Albán in Oaxaca state in Mexico was built around 275 BC, at a different orientation to other structures in the complex. Its steps are aligned perpendicular to the rising of Capella at that time, so that a person looking out a doorway on the building would have faced it directly. Capella is significant as its heliacal rising took place within a day of the Sun passing directly overhead over Monte Albán.
Multiple status
Professor William Wallace Campbell of the Lick Observatory announced that Capella was binary in 1899, based on spectroscopic observations—he noted on photographic plates taken from August 1896 to February 1897 that a second spectrum appeared superimposed over the first, and that there was a doppler shift to violet in September and October and to red in November and February—showing that the components were moving toward and away from the Earth (and hence orbiting each other). Almost simultaneously, British astronomer Hugh Newall had observed its composite spectrum with a four prism spectroscope attached to a telescope at Cambridge in July 1899, concluding that it was a binary star system.
Many observers tried to discern the component stars without success. Known as "The Interferometrist's Friend", it was first resolved interferometrically in 1919 by John Anderson and Francis Pease at Mount Wilson Observatory, who published an orbit in 1920 based on their observations. This was the first interferometric measurement of any object outside the Solar System. A high-precision orbit was published in 1994 based on observations by the Mark III Stellar Interferometer, again at Mount Wilson Observatory. Capella also became the first astronomical object to be imaged by a separate element optical interferometer when it was imaged by the Cambridge Optical Aperture Synthesis Telescope in September 1995.
In 1914, Finnish astronomer Ragnar Furuhjelm observed that the spectroscopic binary had a faint companion star, which, as its proper motion was similar to that of the spectroscopic binary, was probably physically bound to it. In February 1936, Carl L. Stearns observed that this companion appeared to be double itself; this was confirmed in September that year by Gerard Kuiper. This pair are designated Capella H and L.
X-ray source
Two Aerobee-Hi rocket flights on September 20, 1962, and March 15, 1963, detected and confirmed an X-ray source in Auriga at RA Dec , identified as Capella. A major milestone in stellar X-ray astronomy happened on April 5, 1974, with the detection of the strongest emission of X-rays up to that time from Capella, measured at more than 10,000 times the x-ray luminosity of the Sun. A rocket flight on that date briefly calibrated its attitude control system when a star sensor pointed the payload axis at Capella. During this period, X-rays in the range 0.2–1.6 keV were detected by an X-ray reflector system co-aligned with the star sensor.
The X-ray luminosity (Lx) of ~1024 W (1031 erg s−1) is four orders of magnitude above the Sun's X-ray luminosity. Capella's X-rays are thought to be primarily from the corona of the most massive star. Capella is ROSAT X-ray source 1RXS J051642.2+460001. The high temperature of Capella's corona as obtained from the first coronal X-ray spectrum of Capella using HEAO 1 would require magnetic confinement, unless it is a free-flowing coronal wind.
Observation
With an average apparent magnitude of +0.08, Capella is the brightest object in the constellation Auriga, the sixth-brightest star in the night sky, the third-brightest in the northern celestial hemisphere (after Arcturus and Vega), and the fourth-brightest visible to the naked eye from the latitude 40°N. It appears to be a rich yellowish-white colour, although the yellow colour is more apparent during daylight observation with a telescope, due to the contrast against the blue sky.
Capella is closer to the north celestial pole than any other first-magnitude star. Its northern declination is such that it is actually invisible south of latitude 44°S—this includes southernmost New Zealand, Argentina and Chile as well as the Falkland Islands. Conversely it is circumpolar north of 44°N: for the whole of the United Kingdom and Canada (except for part of Southern Ontario), most of Europe, and the northernmost fringes of the contiguous United States, the star never sets. Capella and Vega are on opposite sides of the pole, at about the same distance from it, such that an imaginary line between the two stars will nearly pass through Polaris. Visible halfway between Orion's Belt and Polaris, Capella is at its highest in the night sky at midnight in early December and is regarded as a prominent star of the northern winter sky.
A few degrees to the southwest of Capella lie three stars, Epsilon Aurigae, Zeta Aurigae and Eta Aurigae, the latter two of which are known as "The Kids", or Haedi. The four form a familiar pattern, or asterism, in the sky.
Distance
Based on an annual parallax shift of 76.20 milliarcseconds (with a margin of error of 0.46 milliarcseconds) as measured by the Hipparcos satellite, this system is estimated to be from Earth, with a margin of error of 0.3 light-year (0.09 parsec). An alternative method to determine the distance is via the orbital parallax, which gives a distance of with a margin of error of only 0.1%. Capella is estimated to have been a little closer to the Solar System in the past, passing within 29 light-years distant around 237,000 years ago. At this range, it would have shone at apparent magnitude −0.82, comparable to Canopus today.
In a 1960 paper, American astronomer Olin J. Eggen concluded that Capella was a member of the Hyades moving group, a group of stars moving in the same direction as the Hyades cluster, after analysing its proper motion and parallax. Members of the group are of a similar age, and those that are around 2.5 times as massive as the Sun have moved off the main sequence after exhausting their core hydrogen reserves and are expanding and cooling into red giants.
Stellar system
There are several stars within a few arcminutes of Capella and some have been listed as companions in various multiple star catalogues. The Washington Double Star Catalog lists components A, B, C, D, E, F, G, H, I, L, M, N, O, P, Q, and R, with A being the naked-eye star. Most are only line-of-sight companions, but the close pair of red dwarfs H and L are at the same distance as the bright component A and moving through space along with it. Capella A is itself a spectroscopic binary with components Aa and Ab, both giant stars. The pair of giants is separated from the pair of red dwarfs by 723".
American astronomer Robert Burnham Jr. described a scale model of the system where Capella A was represented by spheres 13 and 7 inches across, separated by ten feet. The red dwarfs were then each 0.7 inch across and they were separated by 420 feet. At this scale, the two pairs are 21 miles apart.
Capella A
Capella A consists of two yellow evolved stars that have been calculated to orbit each other every 104.02128 ± 0.00016 days, with a semimajor axis of 111.11 ± 0.10 million km (0.74272 ± 0.00069 AU), roughly the distance between Venus and the Sun. The pair is not an eclipsing binary—that is, as seen from Earth, neither star passes in front of the other. The orbit is known extremely accurately and can be used to derive an orbital parallax with far better precision than the one measured directly. The stars are not near enough to each other for the Roche lobe of either star to have been filled and any significant mass transfer to have taken place, even during the red giant stage of the primary star.
Modern convention designates the more luminous cooler star as component Aa and its spectral type has been usually measured between G2 and K0. The hotter secondary Ab has been given various spectral types of late (cooler) F or early (warmer) G. The MK spectral types of the two stars have been measured a number of times, and they are both consistently assigned a luminosity class of III indicating a giant star. The composite spectrum appears to be dominated by the primary star due to its sharper absorption lines; the lines from the secondary are broadened and blurred by its rapid rotation. The composite spectral class is given as approximately G3III, but with a specific mention of features due to a cooler component. The most recent specific published types are K0III and G1III, although older values are still widely quoted such as G5IIIe + G0III from the Bright Star Catalogue or G8III + G0III by Eggen. Where the context is clear, these two components have been referred to as A and B.
The individual apparent magnitudes of the two component stars cannot be directly measured, but their relative brightness has been measured at various wavelengths. They have very nearly equal brightness in the visible light spectrum, with the hotter secondary component generally being found to be a few tenths of a magnitude brighter. A 2016 measurement gives the magnitude difference between the two stars at a wavelength of 700 nm as 0.00 ± 0.1.
The physical properties of the two stars can be determined with high accuracy. The masses are derived directly from the orbital solution, with Aa being and Ab being . Their angular radii have been directly measured; in combination with the very accurate distance, this gives and for Aa and Ab, respectively. Their surface temperatures can be calculated by comparison of observed and synthetic spectra, direct measurement of their angular diameters and brightnesses, calibration against their observed colour indices, and disentangling of high resolution spectra. Weighted averages of these four methods give 4,970 ± 50 K for Aa and 5,730 ± 60 for Ab. Their bolometric luminosities are most accurately derived from their apparent magnitudes and bolometric corrections, but are confirmed by calculation from the temperatures and radii of the stars. Aa is 78.7 ± 4.2 times as luminous as the Sun and Ab 72.7 ± 3.6 times as luminous, so the star defined as the primary component is the more luminous when all wavelengths are considered but very slightly less bright at visual wavelengths.
Estimated to be 590 to 650 million years old, the stars were probably at the hot end of spectral class A during their main-sequence lifetime, similar to Vega. They have now exhausted their core hydrogen and evolved off the main sequence, their outer layers expanding and cooling. Despite the giant luminosity class, the secondary component is very clearly within the Hertzsprung gap on the Hertzsprung–Russell diagram, still expanding and cooling towards the red giant branch, making it a subgiant in evolutionary terms. The more massive primary has already passed through this stage, when it reached a maximum radius of 36 to 38 times that of the Sun. It is now a red clump star which is fusing helium to carbon and oxygen in its core, a process that has not yet begun for the less massive star. Detailed analysis shows that it is nearing the end of this stage and starting to expand again which will lead it to the asymptotic giant branch. Isotope abundances and spin rates confirm this evolutionary difference between the two stars. Heavy element abundances are broadly comparable to those of the Sun and the overall metallicity is slightly less than the Sun's.
The rotational period of each star can be measured by observing periodic variations in the doppler shifts of their spectral lines. The absolute rotational velocities of the two stars are known from their inclinations, rotation periods, and sizes, but the projected equatorial rotational velocities measured using doppler broadening of spectral lines are a standard measure and these are generally quoted. Capella Aa has a projected rotational velocity of 4.1 ± 0.4 km per second, taking 104 ± 3 days to complete one rotation, while Capella Ab spins much more rapidly at 35.0 ± 0.5 km per second, completing a full rotation in only 8.5 ± 0.2 days. Rotational braking occurs in all stars when they expand into giants, and binary stars are also tidally braked. Capella Aa has slowed until it is rotationally locked to the orbital period, although theory predicts that it should still be rotating more quickly from a starting point of a rapidly-spinning main sequence A star.
Capella has long been suspected to be slightly variable. Its amplitude of about 0.1 magnitudes means that it may at times be brighter or fainter than Rigel, Betelgeuse and Vega, which are also variable. The system has been classified as an RS Canum Venaticorum variable, a class of binary stars with active chromospheres that cause huge starspots, but it is still only listed as a suspected variable in the General Catalogue of Variable Stars. Unusually for RS CVn systems, the hotter star, Capella Ab, has the more active atmosphere because it is located in the Hertzsprung gap—a stage where it is changing its angular momentum and deepening its convection zone.
The active atmospheres and closeness of these stars means that they are among the brightest X-ray sources in the sky. However the X-ray emission is due to stable coronal structures and not eruptive flaring activity. Coronal loops larger than the Sun and with temperatures of several million kelvin are likely to be responsible for the majority of the X-rays.
Capella HL
The seventh companion published for Capella, component H, is physically associated with the bright primary star. It is a red dwarf separated from the pair of G-type giants by a distance of around . It has its own close companion, an even fainter red dwarf that was 1.8″ away when it was discovered in 1935. It is component L in double star catalogues. In 2015 the separation had increased to 3.5″, which was sufficient to allow tentative orbital parameters to be derived, 80 years after its discovery. The Gliese-Jahreiss Catalogue of nearby stars designates the binary system as GJ 195. The two components are then referred to individually as GJ 195 A and B.
The two stars are reported to have a 3.5-visual-magnitude difference (2.3 mag in the passband of the Gaia spacecraft) although the difference is much smaller at infrared wavelengths. This is unexpected and may indicate further unseen companions.
The mass of the stars can, in principle, be determined from the orbital motion, but uncertainties in the orbit have led to widely varying results. In 1975, an eccentric 388-year orbit gave masses of and . A smaller near-circular orbit published in 2015 had a 300-year orbit, benefitting from mass constraints of and , respectively, for GJ 195 A and B, based on their infrared magnitudes.
Visual companions
Six visual companions to Capella were discovered before Capella H and are generally known only as Capella B through G. None are thought to be physically associated with Capella, although all appear closer in the sky than the HL pair.
Component F is also known as TYC 3358-3142-1. It is listed with a spectral type of K although it is included in a catalogue of OB stars as a distant luminous star.
Component G is BD+45 1076, with a spectral type of F0, at a distance of . It is identified as a variable member of the Guide Star Catalogue from Chandra observations although it is not known what type of variability. It is known to be an X-ray source with an active corona.
Several other stars have also been catalogued as companions to Capella. Components I, Q and R are 13th-magnitude stars at distances of 92″, 133″ and 134″. V538 Aurigae and its close companion HD 233153 are red dwarfs ten degrees away from Capella; they have very similar space motions but the small difference makes it possible that this is just a coincidence. Two faint stars have been discovered by speckle imaging in the Capella HL field, around 10″ distant from that pair. These have been catalogued as Capella O and P. It is not known whether they are physically associated with the red dwarf binary.
Etymology and culture
Capella traditionally marks the left shoulder of the constellation's eponymous charioteer, or, according to the 2nd-century astronomer Ptolemy's Almagest, the goat that the charioteer is carrying. In Bayer's 1603 work Uranometria, Capella marks the charioteer's back. The three Haedi had been identified as a separate constellation by Pliny the Elder and Manilius, and were called Capra, Caper, or Hircus, all of which relate to its status as the "goat star". Ptolemy merged the Charioteer and the Goats in the 2nd-century Almagest.
In Greek mythology, the star represented the goat Amalthea that suckled Zeus. It was this goat whose horn, after accidentally being broken off by Zeus, was transformed into the cornucopia, or "horn of plenty", which would be filled with whatever its owner desired. Though most often associated with Amalthea, Capella has sometimes been associated with Amalthea's owner, a nymph. The myth of the nymph says that the goat's hideous appearance, resembling a Gorgon, was partially responsible for the Titans' defeat, after Zeus skinned the goat and wore it as his aegis.
In medieval accounts, it bore the uncommon name Alhajoth (also spelled Alhaior, Althaiot, Alhaiset, Alhatod, Alhojet, Alanac, Alanat, Alioc), which (especially the last) may be a corruption of its Arabic name, , al-cayyūq. cAyyūq has no clear significance in Arabic, but may be an Arabized form of the Greek αίξ aiks "goat"; cf. the modern Greek Αίγα Aiga, the feminine of goat. To the Bedouin of the Negev and Sinai, Capella al-'Ayyūq ath-Thurayyā "Capella of the Pleiades", from its role as pointing out the position of that asterism. Another name in Arabic was Al-Rākib "the driver", a translation of the Greek.
To the ancient Balts, Capella was known as Perkūno Ožka "Thunder's Goat", or Tikutis. Conversely in Slavic Macedonian folklore, Capella was Jastreb "the hawk", flying high above and ready to pounce on Mother Hen (the Pleiades) and the Rooster (Nath).
Astrologically, Capella portends civic and military honors and wealth. In the Middle Ages, it was considered a Behenian fixed star, with the stone sapphire and the plants horehound, mint, mugwort and mandrake as attributes. Cornelius Agrippa listed its kabbalistic sign with the name Hircus (Latin for goat).
In Hindu mythology, Capella was seen as the heart of Brahma, Brahma Hṛdaya. In traditional Chinese astronomy, Capella was part of the asterism (; English: Five Chariots), which consisted of Capella together with Beta Aurigae, Theta Aurigae and Iota Aurigae, as well as Beta Tauri. Since it was the second star in this asterism, it has the Chinese name (; English: Second of the Five Chariots).
In Quechua it was known as Colça; the Incas held the star in high regard. The Hawaiians saw Capella as part of an asterism Ke ka o Makali'i ("The canoe bailer of Makali'i") that helped them navigate at sea. Called Hoku-lei "star wreath", it formed this asterism with Procyon, Sirius, Castor and Pollux. In Tahitian folklore, Capella was Tahi-ari'i, the wife of Fa'a-nui (Auriga) and mother of prince Ta'urua (Venus) who sails his canoe across the sky. In Inuit astronomy, Capella, along with Menkalinan (Beta Aurigae), Pollux (Beta Geminorum) and Castor (Alpha Geminorum), formed a constellation Quturjuuk, "collar-bones", the two pairs of stars denoting a bone each. Used for navigation and time-keeping at night, the constellation was recognised from Alaska to western Greenland. The Gwich'in saw Capella and Menkalinan has forming shreets'ą įį vidzee, the right ear of the large circumpolar constellation Yahdii, which covered much of the night sky, and whose orientation facilitated navigation and timekeeping.
In Australian Aboriginal mythology for the Boorong people of Victoria, Capella was Purra, the kangaroo, pursued and killed by the nearby Gemini twins, Yurree (Castor) and Wanjel (Pollux). The Wardaman people of northern Australia knew the star as Yagalal, a ceremonial fish scale, related to Guwamba the barramundi (Aldebaran).
Namesakes
Capella, a lunar crater to the north of the Mare Nectaris, not named after the star
and USNS Capella (T-AKR-293), both U.S. Navy ships
Mazda Capella, a model of automobile manufactured by Mazda
| Physical sciences | Notable stars | null |
1264088 | https://en.wikipedia.org/wiki/Hematopoietic%20stem%20cell | Hematopoietic stem cell | Hematopoietic stem cells (HSCs) are the stem cells that give rise to other blood cells. This process is called haematopoiesis. In vertebrates, the first definitive HSCs arise from the ventral endothelial wall of the embryonic aorta within the (midgestational) aorta-gonad-mesonephros region, through a process known as endothelial-to-hematopoietic transition. In adults, haematopoiesis occurs in the red bone marrow, in the core of most bones. The red bone marrow is derived from the layer of the embryo called the mesoderm.
Haematopoiesis is the process by which all mature blood cells are produced. It must balance enormous production needs (the average person produces more than 500 billion blood cells every day) with the need to regulate the number of each blood cell type in the circulation. In vertebrates, the vast majority of hematopoiesis occurs in the bone marrow and is derived from a limited number of hematopoietic stem cells that are multipotent and capable of extensive self-renewal.
Hematopoietic stem cells give rise to different types of blood cells, in lines called myeloid and lymphoid. Myeloid and lymphoid lineages both are involved in dendritic cell formation. Myeloid cells include monocytes, macrophages, neutrophils, basophils, eosinophils, erythrocytes, and megakaryocytes to platelets. Lymphoid cells include T cells, B cells, natural killer cells, and innate lymphoid cells.
The definition of hematopoietic stem cell has developed since they were first discovered in 1961. The hematopoietic tissue contains cells with long-term and short-term regeneration capacities and committed multipotent, oligopotent, and unipotent progenitors. Hematopoietic stem cells constitute 1:10,000 of cells in myeloid tissue.
HSC transplants are used in the treatment of cancers and other immune system disorders due to their regenerative properties.
Structure
They are round, non-adherent, with a rounded nucleus and low cytoplasm-to-nucleus ratio. In shape, hematopoietic stem cells resemble lymphocytes.
Location
The very first hematopoietic stem cells during (mouse and human) embryonic development are found in aorta-gonad-mesonephros region and the vitelline and umbilical arteries. Slightly later, HSCs are also found in the placenta, yolk sac, embryonic head, and fetal liver.
Stem and progenitor cells can be taken from the pelvis, at the iliac crest, using a needle and syringe. The cells can be removed as liquid (to perform a smear to look at the cell morphology) or they can be removed via a core biopsy (to maintain the architecture or relationship of the cells to each other and to the bone).
Subtypes
A colony-forming unit is a subtype of HSC. (This sense of the term is different from colony-forming units of microbes, which is a cell counting unit.) There are various kinds of HSC colony-forming units:
Colony-forming unit–granulocyte-erythrocyte-monocyte-megakaryocyte (CFU-GEMM)
Colony-forming unit–lymphocyte (CFU-L)
Colony-forming unit–erythrocyte (CFU-E)
Colony-forming unit–granulocyte-macrophage (CFU-GM)
Colony-forming unit–megakaryocyte (CFU-Meg)
Colony-forming unit–basophil (CFU-Baso)
Colony-forming unit–eosinophil (CFU-Eos)
The above CFUs are based on the lineage. Another CFU, the colony-forming unit–spleen (CFU-S), was the basis of an in vivo clonal colony formation, which depends on the ability of infused bone marrow cells to give rise to clones of maturing hematopoietic cells in the spleens of irradiated mice after 8 to 12 days. It was used extensively in early studies, but is now considered to measure more mature progenitor or transit-amplifying cells rather than stem cells.
Isolating stem cells
Since hematopoietic stem cells cannot be isolated as a pure population, it is not possible to identify them in a microscope. Hematopoietic stem cells can be identified or isolated by the use of flow cytometry where the combination of several different cell surface markers (particularly CD34) are used to separate the rare hematopoietic stem cells from the surrounding blood cells. Hematopoietic stem cells lack expression of mature blood cell markers and are thus called Lin-. Lack of expression of lineage markers is used in combination with detection of several positive cell-surface markers to isolate hematopoietic stem cells. In addition, hematopoietic stem cells are characterised by their small size and low staining with vital dyes such as rhodamine 123 (rhodamine lo) or Hoechst 33342 (side population).
Function
Haematopoiesis
Hematopoietic stem cells are essential to haematopoiesis, the formation of the cells within blood. Hematopoietic stem cells can replenish all blood cell types (i.e., are multipotent) and self-renew. A small number of hematopoietic stem cells can expand to generate a very large number of daughter hematopoietic stem cells. This phenomenon is used in bone marrow transplantation, when a small number of hematopoietic stem cells reconstitute the hematopoietic system. This process indicates that, subsequent to bone marrow transplantation, symmetrical cell divisions into two daughter hematopoietic stem cells must occur.
Stem cell self-renewal is thought to occur in the stem cell niche in the bone marrow, and it is reasonable to assume that key signals present in this niche will be important in self-renewal. There is much interest in the environmental and molecular requirements for HSC self-renewal, as understanding the ability of HSC to replenish themselves will eventually allow the generation of expanded populations of HSC in vitro that can be used therapeutically.
Quiescence
Hematopoietic stem cells, like all adult stem cells, mostly exist in a state of quiescence, or reversible growth arrest. The altered metabolism of quiescent HSCs helps the cells survive for extended periods of time in the hypoxic bone marrow environment. When provoked by cell death or damage, Hematopoietic stem cells exit quiescence and begin actively dividing again. The transition from dormancy to propagation and back is regulated by the MEK/ERK pathway and PI3K/AKT/mTOR pathway. Dysregulation of these transitions can lead to stem cell exhaustion, or the gradual loss of active Hematopoietic stem cells in the blood system.
Mobility
Hematopoietic stem cells have a higher potential than other immature blood cells to pass the bone marrow barrier, and, thus, may travel in the blood from the bone marrow in one bone to another bone. If they settle in the thymus, they may develop into T cells. In the case of fetuses and other extramedullary hematopoiesis. Hematopoietic stem cells may also settle in the liver or spleen and develop.
This enables Hematopoietic stem cells to be harvested directly from the blood.
Clinical significance
Transplant
Hematopoietic stem cell transplantation (HSCT) is the transplantation of multipotent hematopoietic stem cells, usually derived from bone marrow, peripheral blood, or umbilical cord blood. It may be autologous (the patient's own stem cells are used), allogeneic (the stem cells come from a donor) or syngeneic (from an identical twin).
It is most often performed for patients with certain cancers of the blood or bone marrow, such as multiple myeloma or leukemia. In these cases, the recipient's immune system is usually destroyed with radiation or chemotherapy before the transplantation. Infection and graft-versus-host disease are major complications of allogeneic HSCT.
In order to harvest stem cells from the circulating peripheral blood, blood donors are injected with a cytokine, such as granulocyte-colony stimulating factor (G-CSF), that induces cells to leave the bone marrow and circulate in the blood vessels.
In mammalian embryology, the first definitive Hematopoietic stem cells are detected in the AGM (aorta-gonad-mesonephros), and then massively expanded in the fetal liver prior to colonising the bone marrow before birth.
Hematopoietic stem cell transplantation remains a dangerous procedure with many possible complications; it is reserved for patients with life-threatening diseases. As survival following the procedure has increased, its use has expanded beyond cancer to autoimmune diseases and hereditary skeletal dysplasias; notably malignant infantile osteopetrosis and mucopolysaccharidosis.
Stem cells can be used to regenerate different types of tissues. HCT is an established as therapy for chronic myeloid leukemia, acute lymphatic leukemia, aplastic anemia, and hemoglobinopathies, in addition to acute myeloid leukemia and primary immune deficiencies. Hematopoietic system regeneration is typically achieved within 2–4 weeks post-chemo- or irradiation therapy and HCT. HSCs are being clinically tested for their use in non-hematopoietic tissue regeneration.
Aging of hematopoietic stem cells
DNA damage
DNA strand breaks accumulate in long term hematopoietic stem cells during aging. This accumulation is associated with a broad attenuation of DNA repair and response pathways that depends on HSC quiescence. Non-homologous end joining (NHEJ) is a pathway that repairs double-strand breaks in DNA. NHEJ is referred to as "non-homologous" because the break ends are directly ligated without the need for a homologous template. The NHEJ pathway depends on several proteins including ligase 4, DNA polymerase mu and NHEJ factor 1 (NHEJ1, also known as Cernunnos or XLF).
DNA ligase 4 (Lig4) has a highly specific role in the repair of double-strand breaks by NHEJ. Lig4 deficiency in the mouse causes a progressive loss of hematopoietic stem cells during aging. Deficiency of lig4 in pluripotent stem cells results in accumulation of DNA double-strand breaks and enhanced apoptosis.
In polymerase mu mutant mice, hematopoietic cell development is defective in several peripheral and bone marrow cell populations with about a 40% decrease in bone marrow cell number that includes several hematopoietic lineages. Expansion potential of hematopoietic progenitor cells is also reduced. These characteristics correlate with reduced ability to repair double-strand breaks in hematopoietic tissue.
Deficiency of NHEJ factor 1 in mice leads to premature aging of hematopoietic stem cells as indicated by several lines of evidence including evidence that long-term repopulation is defective and worsens over time. Using a human induced pluripotent stem cell model of NHEJ1 deficiency, it was shown that NHEJ1 has an important role in promoting survival of the primitive hematopoietic progenitors. These NHEJ1 deficient cells possess a weak NHEJ1-mediated repair capacity that is apparently incapable of coping with DNA damages induced by physiological stress, normal metabolism, and ionizing radiation.
The sensitivity of hematopoietic stem cells to Lig4, DNA polymerase mu and NHEJ1 deficiency suggests that NHEJ is a key determinant of the ability of stem cells to maintain themselves against physiological stress over time. Rossi et al. found that endogenous DNA damage accumulates with age even in wild type Hematopoietic stem cells, and suggested that DNA damage accrual may be an important physiological mechanism of stem cell aging.
Loss of clonal diversity
A study shows the clonal diversity of hematopoietic stem cells gets drastically reduced around age 70 , substantiating a novel theory of ageing which could enable healthy aging. Of note, the shift in clonal diversity during aging was previously reported in 2008 for the murine system by the Christa Muller-Sieburg laboratory in San Diego, California.
Research
Behavior in culture
A cobblestone area-forming cell (CAFC) assay is a cell culture-based empirical assay. When plated onto a confluent culture of stromal feeder layer, a fraction of hematopoietic stem cells creep between the gaps (even though the stromal cells are touching each other) and eventually settle between the stromal cells and the substratum (here the dish surface) or trapped in the cellular processes between the stromal cells. Emperipolesis is the in vivo phenomenon in which one cell is completely engulfed into another (e.g. thymocytes into thymic nurse cells); on the other hand, when in vitro, lymphoid lineage cells creep beneath nurse-like cells, the process is called pseudoemperipolesis. This similar phenomenon is more commonly known in the HSC field by the cell culture terminology cobble stone area-forming cells (CAFC), which means areas or clusters of cells look dull cobblestone-like under phase contrast microscopy, compared to the other hematopoietic stem cells, which are refractile. This happens because the cells that are floating loosely on top of the stromal cells are spherical and thus refractile. However, the cells that creep beneath the stromal cells are flattened and, thus, not refractile. The mechanism of pseudoemperipolesis is only recently coming to light. It may be mediated by interaction through CXCR4 (CD184) the receptor for CXC Chemokines (e.g., SDF1) and α4β1 integrins.
Repopulation kinetics
Hematopoietic stem cells (HSC) cannot be easily observed directly, and, therefore, their behaviors need to be inferred indirectly. Clonal studies are likely the closest technique for single cell in vivo studies of HSC. Here, sophisticated experimental and statistical methods are used to ascertain that, with a high probability, a single HSC is contained in a transplant administered to a lethally irradiated host. The clonal expansion of this stem cell can then be observed over time by monitoring the percent donor-type cells in blood as the host is reconstituted. The resulting time series is defined as the repopulation kinetic of the HSC.
The reconstitution kinetics are very heterogeneous. However, using symbolic dynamics, one can show that they fall into a limited number of classes. To prove this, several hundred experimental repopulation kinetics from clonal Thy-1lo SCA-1+ lin−(B220, CD4, CD8, Gr-1, Mac-1 and Ter-119) c-kit+ HSC were translated into symbolic sequences by assigning the symbols "+", "-", "~" whenever two successive measurements of the percent donor-type cells have a positive, negative, or unchanged slope, respectively. By using the Hamming distance, the repopulation patterns were subjected to cluster analysis yielding 16 distinct groups of kinetics. To finish the empirical proof, the Laplace add-one approach was used to determine that the probability of finding kinetics not contained in these 16 groups is very small. By corollary, this result shows that the hematopoietic stem cell compartment is also heterogeneous by dynamical criteria.
It was originally believed that all hematopoietic stem cells were alike in their self-renewal and differentiation abilities. This view was first challenged by the 2002 discovery by the Muller-Sieburg group in San Diego, who illustrated that different stem cells can show distinct repopulation patterns that are epigenetically predetermined intrinsic properties of clonal Thy-1lo Sca-1+ lin− c-kit+ HSC. The results of these clonal studies led to the notion of lineage bias. Using the ratio of lymphoid (L) to myeloid (M) cells in blood as a quantitative marker, the stem cell compartment can be split into three categories of HSC. Balanced (Bala) hematopoietic stem cells repopulate peripheral white blood cells in the same ratio of myeloid to lymphoid cells as seen in unmanipulated mice (on average about 15% myeloid and 85% lymphoid cells, or 3 ≤ ρ ≤ 10). Myeloid-biased (My-bi) hematopoietic stem cells give rise to very few lymphocytes resulting in ratios 0 < ρ < 3, while lymphoid-biased (Ly-bi) hematopoietic stem cells generate very few myeloid cells, which results in lymphoid-to-myeloid ratios of ρ > 10. All three types are normal types of HSC, and they do not represent stages of differentiation. Rather, these are three classes of HSC, each with an epigenetically fixed differentiation program. These studies also showed that lineage bias is not stochastically regulated or dependent on differences in environmental influence. My-bi HSC self-renew longer than balanced or Ly-bi HSC. The myeloid bias results from reduced responsiveness to the lymphopoetin interleukin 7 (IL-7).
Subsequently, other groups confirmed and highlighted the original findings. For example, the Eaves group confirmed in 2007 that repopulation kinetics, long-term self-renewal capacity, and My-bi and Ly-bi are stably inherited intrinsic HSC properties. In 2010, the Goodell group provided additional insights about the molecular basis of lineage bias in side population (SP) SCA-1+ lin− c-kit+ HSC. As previously shown for IL-7 signaling, it was found that a member of the transforming growth factor family (TGF-beta) induces and inhibits the proliferation of My-bi and Ly-bi HSC, respectively.
Etymology
From Greek haimato-, combining form of haima 'blood', and from the Latinized form of Greek poietikos 'capable of making, creative, productive', from poiein 'to make, create'.
| Biology and health sciences | Cell processes | Biology |
1264197 | https://en.wikipedia.org/wiki/Swather | Swather | A swather (North America), or windrower (Australia and rest of world), is a farm implement that cuts hay or small grain crops and forms them into a windrow for drying.
They may be self-propelled with an engine, or drawn by a tractor and power take-off powered. A swather uses a reciprocating sickle bar or rotating discs to sever the crop stems. The reel helps cut crop fall neatly onto a canvas or auger conveyor which deposits it into a windrow with stems aligned and supported above the ground by the stubble.
A swather does the same task for hay crops as hand scything, cradling and swathing, or mowing and raking. Horizontal rollers behind the cutters may be used to crimp or condition the stems of hay crops to decrease drying time.
For grains, as combines replaced threshing machines, the swather introduced an optional step in the harvesting process to provide for the drying time that binding formerly afforded. Swathing is still more common in the northern United States and Canada than regions with longer growing seasons where standing grain crops can be harvested directly by combines. Some modern crop varieties capable of rapid maturity have reduced the need for swathing grains even in the north.
As well as accelerating drying of the ripe grain, windrowing the whole of the growing crop provides for a consistent ripening and dehydration of stalk and green weeds to assist in effective post threshing winnowing and separation of the grain and other material. Alternatively, chemical desiccation of weedy or irregularly ripe standing crops with glyphosate, paraquat or diquat has been used to enable direct combining.
A swather is the mascot of sports teams at Hesston High School in Hesston, Kansas. Hesston is home to AGCO Corporation swather and combine harvester manufacturing plants.
| Technology | Farm and garden machinery | null |
1264212 | https://en.wikipedia.org/wiki/Elasmosaurus | Elasmosaurus | Elasmosaurus () is a genus of plesiosaur that lived in North America during the Campanian stage of the Late Cretaceous period, about 80.5million years ago. The first specimen was discovered in 1867 near Fort Wallace, Kansas, US, and was sent to the American paleontologist Edward Drinker Cope, who named it E.platyurus in 1868. The generic name means "thin-plate reptile", and the specific name means "flat-tailed". Cope originally reconstructed the skeleton of Elasmosaurus with the skull at the end of the tail, an error which was made light of by the paleontologist Othniel Charles Marsh, and became part of their "Bone Wars" rivalry. Only one incomplete Elasmosaurus skeleton is definitely known, consisting of a fragmentary skull, the spine, and now lost pectoral and pelvic girdles, and a single species is recognized today; other species are now considered invalid or have been moved to other genera.
Measuring in length, Elasmosaurus would have had a streamlined body with paddle-like limbs, a short tail, a small head, and an extremely long neck. The neck alone was around long. Along with its relative Albertonectes, it was one of the longest-necked animals to have lived, with the second largest number of neck vertebrae known, 72, 4 fewer than Albertonectes. The skull would have been slender and triangular, with large, fang-like teeth at the front, and smaller teeth towards the back. It had six teeth in each premaxilla of the upper jaw, and may have had 14 teeth in the maxilla and 19 in the dentary of the lower jaw. Most of the neck vertebrae were compressed sideways, and bore a longitudinal crest or keel along the sides.
The family Elasmosauridae was based on the genus Elasmosaurus, the first recognized member of this group of long-necked plesiosaurs. Elasmosaurids were well adapted for aquatic life, and used their flippers for swimming. Contrary to earlier depictions, their necks were not very flexible, and could not be held high above the water surface. It is unknown what their long necks were used for, but they may have had a function in feeding. Elasmosaurids probably ate small fish and marine invertebrates, seizing them with their long teeth, and may have used gastroliths (stomach stones) to help digest their food. Elasmosaurus is known from the Pierre Shale formation, which represents marine deposits from the Western Interior Seaway.
History of study
In early 1867, the American army surgeon Theophilus Hunt Turner and the army scout William Comstock explored the rocks around Fort Wallace, Kansas, where they were stationed during the construction of the Union Pacific Railroad. Approximately northeast of Fort Wallace, near McAllaster, Turner discovered the bones of a large fossil reptile in a ravine in the Pierre Shale formation, and though he had no paleontological experience, he recognized the remains as belonging to an "extinct monster". In June, Turner gave three fossil vertebrae to the American scientist John LeConte, a member of the railway survey, to take back east to be identified. In December, LeConte delivered some of the vertebrae to the American paleontologist Edward Drinker Cope at the Academy of Natural Sciences of Philadelphia (ANSP, known since 2011 as the Academy of Natural Sciences of Drexel University). Recognizing them as the remains of a plesiosaur, larger than any he had seen in Europe, Cope wrote to Turner asking him to deliver the rest of the specimen, at the ANSP's expense.
In December 1867 Turner and others from Fort Wallace returned to the site and recovered much of the vertebral column, as well as concretions that contained other bones; the material had a combined weight of . The fossils were dug or pried out of the relatively soft shale with picks and shovels, loaded on a horse-drawn wagon, and transported back to Fort Wallace. Cope sent instructions on how to pack the bones, which were thereafter sent in hay-padded crates on a military wagon east to the railroad, which had not yet reached the fort. The specimen arrived in Philadelphia by rail in March 1868, whereafter Cope examined it hurriedly; he reported on it at the March ANSP meeting, during which he named it Elasmosaurus platyurus. The generic name Elasmosaurus means "thin-plate reptile", in reference to the "plate" bones of the sternal and pelvic regions, and the specific name platyurus means "flat-tailed", in reference to the compressed "tail" (actually the neck) and laminae of the vertebrae there.
Cope requested that Turner search for more parts of the Elasmosaurus specimen, and was sent more fossils during August or September 1868. The ANSP thanked Turner for his "very valuable gift" at their meeting in December 1868, and Turner visited the museum during spring, at a time when Cope was absent. Turner died unexpectedly at Fort Wallace on July27, 1869, without seeing the completion of the work he began, but Cope continued to write him, unaware of his death until 1870. The circumstances around Turner's discovery of the type specimen were not covered in Cope's report, and remained unknown until Turner's letters were published in 1987. Elasmosaurus was the first major fossil discovery in Kansas (and the largest from there at the time), and marked the beginning of a fossil collecting rush that sent thousands of fossils from Kansas to prominent museums on the American east coast. Elasmosaurus was one of few plesiosaurs known from the New World at the time, and the first recognized member of the long-necked family of plesiosaurs, the Elasmosauridae.
In 1869 Cope scientifically described and figured Elasmosaurus, and the preprint version of the manuscript contained a reconstruction of the skeleton which he had earlier presented during his report at an ANSP meeting in September 1868. The reconstruction showed Elasmosaurus with a short neck and a long tail, unlike other plesiosaurs, and Cope was also unsure whether it had hind limbs. At an ANSP meeting a year and a half later, in March 1870, the American paleontologist Joseph Leidy (Cope's mentor) noted that Cope's reconstruction of Elasmosaurus showed the skull at the wrong end of the vertebral column, at the end of the tail instead of the neck. Cope had apparently concluded that the tail vertebrae belonged to the neck, since the jaws had been found at that end of the skeleton, even though the opposite end terminated in the axis and atlas bones that are found in the neck. Leidy also concluded that Elasmosaurus was identical to Discosaurus, a plesiosaur he had named in 1851.
To hide his mistake, Cope attempted to recall all copies of the preprint article, and printed a corrected version with a new skeletal reconstruction that placed the head on the neck (though it reversed the orientation of the individual vertebrae) and different wording in 1870. In a reply to Leidy, Cope claimed that he had been misled by the fact that Leidy had arranged the vertebrae of Cimoliasaurus in the reverse order in his 1851 description of that genus, and pointed out that his reconstruction had been corrected. Cope also rejected the idea that Elasmosaurus and Discosaurus were identical, and noted that the latter and Cimoliasaurus did not have any distinguishing features. Though Cope had tried to destroy the preprints, one copy came to the attention of the American paleontologist Othniel Charles Marsh, who made light of the mistake. This led to antagonism between Cope, who was embarrassed by the mistake, and Marsh, who brought up the mistake repeatedly for decades. Marsh returned to the issue during their controversy in the New York Herald in the 1890s (Marsh claimed he had pointed out the error to Cope immediately), when their dispute gained widespread public attention. The argument was part of the "Bone Wars" rivalry between the two, and is well known in the history of paleontology.
Because of Cope's reputation as a brilliant paleontologist, it has been questioned why he would make such an obvious anatomical error. It has been suggested that, as a unique specimen in 1868, the original Elasmosaurus may have been hard to interpret based on the knowledge available at the time. Also, Cope initially thought it consisted of two specimens of different animalsin an 1868 letter to LeConte, Cope had referred to the supposed "smaller specimen" as Discosaurus carinatus. Cope was only in his late twenties and not formally trained in paleontology, and may have been influenced by Leidy's mistake of reversing the vertebral column of Cimoliasaurus. In 2002 the American art historian Jane P. Davidson noted that the fact that other scientists early on had pointed out Leidy's error argues against this explanation, adding that Cope was not convinced he had made a mistake. Plesiosaur anatomy was sufficiently well known at the time that Cope should not have made the mistake, according to Davidson. Cope did little work on the specimen since his 1870 description, and it was kept in storage for nearly 30 years. It was only redescribed in detail in 2005 by the German paleontologist Sven Sachs.
Known and possible fossil elements
Today, the incomplete holotype specimen, cataloged as ANSP 10081, is the only definite specimen of Elasmosaurus. It was long exhibited, but is now stored in a cabinet with other assigned fragments. The specimen consists of the premaxillae, part of the hind-section of the right maxilla, two maxilla fragments with teeth, the front part of the dentaries, three more jaw fragments, two cranial fragments of indeterminable identity, 72 cervical vertebrae of the neck, including the atlas and axis, 3pectoral vertebrae, 6back vertebrae, 4sacral vertebrae, 18 tail vertebrae, as well as rib fragments. In 2013 an incomplete neck vertebra centrum of the holotype that had been mentioned by Cope but thought to have been lost was rediscovered in storage by Sachs, and the neck vertebra count was revised from 71 to 72. The neck vertebrae have been taphonomically distorted (changes occurring during decay and fossilization), with some parts being unnaturally compressed or displaced. In 1986 a three-dimensional reconstruction of the holotype skeleton was completed and is now displayed at the ANSP. This cast was later copied by the company Triebold Paleontology Incorporated, and replicas were provided to other museums. The replica at the Fort Wallace Museum measures about in length.
Though Cope described and figured the pectoral and pelvic girdles of Elasmosaurus in 1869 and 1875, these elements were noted as missing from the collection by the American paleontologist Samuel Wendell Williston in 1906. Cope had loaned these elements to the English sculptor Benjamin Waterhouse Hawkins to help prepare them out of their surrounding concretions. At the time, Hawkins was working on a "Paleozoic Museum" in New York's Central Park, where a reconstruction of Elasmosaurus was to appear, an American equivalent to his life-sized Crystal Palace Dinosaurs in London. In May 1871 much of the exhibit material in Hawkins' workshop was destroyed by vandals for unclear reasons and their fragments buried; it is possible that the girdle elements of Elasmosaurus were at the workshop and were likewise destroyed. Nothing was subsequently mentioned about their loss by Hawkins or Cope. In 2018, Davidson and Everhart documented the events leading up to the disappearance of these fossils, and suggested that a photo and drawing of Waterhouse's workshop from 1869 appear to show concretions on the floor that may have been the unprepared girdles of Elasmosaurus. They also noted that conceptual sketches of the Palaeozoic Museum show that the model Elasmosaurus was originally envisioned with a long "tail", though later updated with a long neck. Davidson and Everhart concluded that the girdle fossils were most likely destroyed in Hawkins' workshop.
Fossils that may have belonged to the holotype were found by the American geologist Benjamin Franklin Mudge in 1871, but have probably been lost since. Additional plesiosaur fossils were recovered near the original locality in 1954, 1991, 1994, and 1998, including back vertebrae, ribs, gastralia (belly ribs), and gastroliths. As none of these elements overlap with those of the holotype specimen, in 2005 the American paleontologist Michael J. Everhart concluded they belonged to the same individual, and that the parts had been separated before burial of the carcass. He also noted that a small stone wedged in the neural canal of one of the tail vertebrae of the holotype may be a gastrolith, based on its polished appearance. In 2007 the Colombian paleontologists Leslie Noè and Marcela Gómez-Pérez expressed doubt that the additional elements belonged to the type specimen, or even to Elasmosaurus, due to lack of evidence. They explained that elements missing from the holotype may have been lost to weathering or simply not collected, and that parts may have been lost or damaged during transportation or preparation. Gastroliths may also not have been recognized as such during collection, since such stones were not reported from a plesiosaur until ten years after.
In 2017 Sachs and Joachim Ladwig suggested that a fragmentary elasmosaurid skeleton from the upper Campanian of Kronsmoor in Schleswig-Holstein, Germany, and housed in the Naturkunde-Museum Bielefeld, may have belonged to Elasmosaurus. Additional parts of the same skeleton are housed at the Institute for Geology of the University of Hamburg, as well as in private collections. Combined, the specimen consists of neck, back and tail vertebrae, phalanges, a tooth, limb elements, 110 gastroliths, and unidentifiable fragments.
Description
Though the only known specimen of Elasmosaurus (holotype specimen ANSP10081) is fragmentary and missing many elements, related elasmosaurids show it would have had a compact, streamlined body, long, paddle-like limbs, a short tail, a proportionately small head, and an extremely long neck. The neck of Elasmosaurus is estimated at in length; thus, Elasmosaurus and its relative Albertonectes were some of the longest-necked animals ever to have lived, with the largest number of neck vertebrae of any known vertebrate animals. In spite of their many neck vertebrae, the necks of elasmosaurids were less than half as long as those of the longest-necked sauropod dinosaurs. Initially, in his 1869 description of Elasmosaurus, Cope estimated the length of the animal by summing up vertebral lengths and estimations of missing parts, resulting in a total length of ; he believed that the living animal would have been slightly larger due to cartilage present between the vertebral bodies, and was estimated at roughly . However, in 1952, the American paleontologist Samuel Welles estimated the body length to have been , a number that was repeated by José Patricio O'Gorman in 2016.
Like other elasmosaurids, Elasmosaurus would have had a slender, triangular skull. The snout was rounded and almost formed a semi-circle when viewed from above, and the premaxillae (which form the front of the upper jaw) bore a low keel at the midline. It is uncertain how many teeth Elasmosaurus had, due to the fragmentary state of the fossils. It probably had six teeth in each premaxilla, and the teeth preserved there were formed like large fangs. The number of premaxillary teeth distinguished Elasmosaurus from primitive plesiosauroids and most other elasmosaurids, which usually had fewer. The two teeth at the front were smaller than the succeeding ones, and were located between the first two teeth in the dentaries of the lower jaws. The known teeth of the front part of the lower jaw were large fangs, and the teeth at the back of the jaws appear to have been smaller. The dentition of elasmosaurids was generally heterodont (irregular throughout the jaws), with the teeth becoming progressively smaller from front to back. The maxillae (largest tooth bearing bone of the upper jaw) of elasmosaurids usually contained 14teeth, whereas the dentaries (the main part of the lower jaws) usually contained 17 to 19. The teeth interlocked, and their tooth crowns were slender and rounded in cross-section. The mandibular symphysis (where the two halves of the lower jaw connected) was well ossified, with no visible suture.
The pectoral and pelvic girdles of the holotype specimen were noted as missing by 1906, but observations about these elements were since made based on the original descriptions and figures from the late 19thcentury. The shoulder blades (scapulae) were fused and met at the midline, bearing no trace of a median bar. The upper processes of the shoulder blades were very broad, and the "necks" of the shoulder blades were long. The pectoral girdle had a long bar at the middle, a supposedly advanced feature thought to be absent from juvenile plesiosaurs. The ischia (a pair of bones that formed part of the pelvis) were joined at the middle, so that a medial bar was present along the length of the pelvis, a feature usually not found in plesiosaurs. Like other elasmosaurids (and plesiosaurs in general), Elasmosaurus would have had large, paddle-like limbs with very long digits. The paddles at the front (the pectoral paddles) were longer than those at the back (the pelvic paddles).
Vertebrae
Unlike those of many other long-necked animals, the individual neck vertebrae were not particularly elongated; rather, the extreme neck length was achieved by a much increased number of vertebrae. Elasmosaurus differed from all other plesiosaurs by having 72 neck vertebrae; more may have been present but were later lost to erosion or after excavation. Only Albertonectes had more neck vertebrae, 76, and the two are the only plesiosaurs with a count higher than 70; more than 60 vertebrae is very derived (or "advanced") for plesiosaurs.
The atlas and axis bone complex, consisting of the first two neck vertebrae and articulated with the back of the skull, was long, low, and horizontally rectangular in side-view. The centra, or "bodies", of these vertebrae were co-ossified in the holotype specimen, which indicates it was an adult. The neural arches of these vertebrae were very thin and rather high, which gave the neural canal (the opening through the middle of the vertebrae) a triangular outline when seen from the back. The lower part of the neural canal was narrow towards the back by the axis, where it was half the breadth of the centrum. It became broader towards the front, where it was almost the same breadth as the centrum of the atlas. The neural arches were also more robust there than in the axis, and the neural canal was higher. The neural spine was low and directed upwards and back. The centra of the atlas and axis were of equal length, and had a quadratic shape in side view. The surface (or facet) where the axis articulated with the next vertebra had an oval outline, and an excavation for the neural canal in the middle of its upper edge. A distinct keel ran along the lower middle of the atlas and axis vertebrae.
Most of the neck vertebrae were compressed sideways, especially at the middle of the neck. A crest (also termed ridge or keel) ran longitudinally along the side of the neck vertebrae (a feature typical of elasmosaurids), visible from the third to the fifty-fifth vertebrae, at the hind part of the neck. This crest was positioned at the middle of the centrum in the front vertebrae, and at the upper half of the centrum from the 19th vertebra and onwards. The crest would have served to anchor the musculature of the neck. The centra differed in shape depending on the position of the vertebrae in the neck; that of the third vertebra was about as long as it was broad, but the centra became longer than broad from the fourth vertebra and onwards. The centra became more elongated at the middle of the neck, but became shorter again at the back of the neck, with the length and breadth being about equal at the 61st vertebra, and those of the hindmost vertebrae being broader than long. The articular surfaces of the vertebrae in the front of the neck were broad oval, and moderately deepened, with rounded, thickened edges, with an excavation (or cavity) at the upper and lower sides. Further back in the front part of the neck, around the 25th vertebra, the lower edge of the articular facets became more concave, and the facet shaped like a quadrate with rounded edges. By the 63rd vertebra, the articular facet was also quadratic in shape with rounded edges, whereas the centra of the hindmost vertebrae had a broad oval outline.
The neural arches of the neck vertebrae were well fused to the centra, leaving no visible sutures, and the neural canal was narrow in the front vertebrae, becoming more prominently developed in the hind vertebrae, where it was as broad as high, and almost circular. The pre-and post-zygapophyses of the neck vertebrae, processes that articulated adjacent vertebrae so they fit together, were of equal length; the former reached entirely over the level of the centrum whereas the latter reached only with their back half. The neural spines of the neck vertebrae appear to have been low, and almost semi-circular by the 20th vertebra. The facets where the neck ribs articulated with the neck vertebrae were placed on the lower sides of the centra, but were only placed higher in the last three vertebrae, reaching around the middle of the sides. The neck ribs were semicircular to quadratic in side view, and were directed rather straight down. The bottom of each neck vertebrae had pairs of nutritive foramina (openings) at the middle, separated by a ridge, which became progressively more prominent and thickened towards the back of the neck.
The vertebrae that transitioned between the neck and back vertebrae in the pectoral region of plesiosaurs, close to the front margin of the forelimb girdle, are often termed pectoral vertebrae. Elasmosaurus had three pectoral vertebrae, which is a common number for elasmosaurids. The rib facets of the pectoral vertebrae were triangular in shape and situated on transverse processes, and the centra bore pairs of nutritive foramina in the middle of the lower sides. The back vertebrae had rib facets level with the neural canal, and the front and back part of the transverse processes here had distinct ridges on their margins. Here the rib facets where placed higher than the transverse processes, separating the two, and were oval to rectangular in outline. The pre-zygapophyses here were shorter than those in the neck and pectoral vertebrae, and only reached above the level of the centrum with the front third of their length. The post-zygapophyses reached over the level of the centrum with the back half of their length. Back vertebrae are not useful for distinguishing between elasmosaurids, since they are not diagnostic at the genus level.
Elasmosaurus had four sacral vertebrae (the fused vertebrae that form the sacrum connected to the pelvis), a number typical of elasmosaurids. The transverse processes here were very short, and the rib facets increased in size from the first to the fourth sacral vertebra. A ridge ran along the top of these vertebrae, and the lower sides of the centra were rounded, and bore pairs of nutritive foramina, separated by low ridges. The first tail (or caudal) vertebra could be distinguished by the preceding sacral vertebra by having smaller rib facets, and by being positioned in the lower half of the centrum. These vertebrae were almost circular in shape, and the first two bore a narrow keel in the middle of the upper side. The rib facets of the tail vertebrae were located on the lower side of the centra, and their oval shape became larger and broader from the third vertebra and onwards, but became smaller from the 14th vertebra. Here, the pre-zygapophyses also reached over the level of the centra for most of their length, while the post-zygapophyses reached over this level by half their length. The lower part of the centra were rounded from the first to the third tail vertebrae, but concave from the fourth to the 18th. The usual number of tail vertebrae in elasmosaurids is 30. Since the last tail-vertebrae of elasmosaurids were fused into a structure similar to the pygostyle of birds, it is possible this supported a tail-fin, but the shape it would have had is unknown.
Formerly assigned species
Following the description of the type species, E. platyurus, a number of other Elasmosaurus species were described by Cope, Williston, and other authors. However, none of these species are still definitely referable to the genus Elasmosaurus today, and most of them either have been moved to genera of their own or are considered dubious names, nomina dubiathat is, with no distinguishing features, and therefore of questionable validity.
Accompanying his 1869 description of E. platyurus, Cope named another species of Elasmosaurus, E.orientalis, based on two back vertebrae from New Jersey. He distinguished E.orientalis from E.platyurus by the more strongly developed processes known as parapophyses on the vertebrae, in which he considered it to approach closer to Cimoliasaurus; however, he still assigned it to Elasmosaurus on account of its large size and angled sides. The first of these vertebrae was used as a doorstop in a tailor's shop, whereas the other was found in a pit by Samuel Lockwood, a superintendent. Cope gave the name orientalis to the new species, on account of it possibly having a more easterly distribution than E.platyurus. Leidy subsequently moved E.orientalis to the now dubious genus Discosaurus in the following year. In 1952 Welles considered the species a nomen dubium, given how fragmentary it was.
In 1869 Cope also published an article about the fossil reptiles of New Jersey, wherein he described E.orientalis as an animal with a "long neck". Yet, in an accompanying illustration Cope showed a short-necked Elasmosaurus confronting a Dryptosaurus (then Laelaps), with a plesiosaur-like Mosasaurus and other animals in the background. According to Davidson, it is uncertain which species of Elasmosaurus is depicted, but if it is E.orientalis, the short neck contradicts Cope's own text, and if E.platyurus, he showed the animal with a short neck after acknowledging this was incorrect. Davidson has suggested that even though Leidy had pointed out Cope's error in 1868, Cope may not have accepted this. In an 1870 reply to Leidy, Cope himself stated that the generic placement of E.orientalis was in doubt, and that he had illustrated it with a short neck due to believing this was the condition of Cimoliasaurus. If more remains showed E.orientalis to have had a long neck like Elasmosaurus, he stated the image may instead represent Cimoliasaurus better.
In the same 1869 publication wherein he named E. platyurus and E.orientalis, Cope assigned an additional species, E.constrictus, based on a partial centrum from a neck vertebra found in the Turonian-aged clay deposits at Steyning, Sussex, in the United Kingdom. It was described by the British paleontologist Richard Owen as Plesiosaurus constrictus in 1850; Owen named the species after the extremely narrow breadth of the vertebra between the pleurapophyses, or the processes that articulate between the ribs. He considered this to be partially an artifact of preservation, but could not understand how the compression affected only the central portion and not the articular ends of the centrum. Cope recognized this as a natural condition, and considered constrictus to be "a species of Elasmosaurus or an ally". In 1962 Welles considered P.constrictus to be a nomen dubium, given its fragmentary nature. Per Ove Persson retained it as valid in 1963, noting the longitudinal ridge on the sides of the centra as an elasmosaurid trait. In 1995 Nathalie Bardet and Pascal Godefroit also recognized it as an elasmosaurid, albeit indeterminate.
Cope discovered another elasmosaurid skeleton in 1876. He named it as a new species, E.serpentinus, in 1877, and differentiated it by the lack of compression in the rear neck vertebrae, the presence of few sessile ribs among the first few dorsals, and the presence of "weak angles" below the front tail vertebrae. Cope had also discovered another large skeleton that bore great resemblance to the known remains of E.orientalis from the black shale of the "Cretaceous bed No.4"; he excavated it with the help of George B. Cledenning and Capt. Nicholas Buesen. In 1943 Welles removed E.serpentinus from Elasmosaurus, and placed it in a new genus, Hydralmosaurus. Subsequently, all Hydralmosaurus specimens were moved to Styxosaurus in 2016, rendering the former a nomen dubium. Williston published a figure of another E. serpentinus specimen in 1914; Elmer Riggs formally described it in 1939. Welles moved this specimen to the new genus and species Alzadasaurus riggsi in 1943. Kenneth Carpenter reassigned it to Thalassomedon haningtoni in 1999; Sachs, Johan Lindgren, and Benjamin Kear noted that the remains represented a juvenile and were significantly distorted, and preferred to retain it as a nomen dubium in 2016.
Subsequently, a series of 19 neck and back vertebrae from the Big Bend region of the Missouripart of the Pierre Shale formationwere found by John H. Charles. Cope, upon receiving the bones at the Academy of Natural Sciences, considered them yet another species of Elasmosaurus. The vertebrae were, according to Cope, the shortest among members of the genus (approaching Cimoliasaurus in this condition), but he still considered them as belonging to Elasmosaurus due to their compressed form. He named it E.intermedius in 1894. However, in his 1906 revision of North American plesiosaurs, Williston regarded the vertebrae as "all more or less mutilated", and found no distinct differences between the remains of E.intermedius and E.platyurus. In 1952 Welles opined that, if E.intermedius was valid, "it must be referred to a pliosaurian genus"; however, he proceeded to label it a nomen dubium in 1962. Three shorter vertebrae found alongside E.intermedius, assigned by Cope to the new genus and species Embaphias circulosus, were also considered by Welles to be a nomen dubium in 1962.
Williston named a number of other new Elasmosaurus species in his 1906 revision. In 1874 he and Mudge discovered a specimen in Plum Creek, Kansas. While he initially assigned it in 1890 to a new species of Cimoliasaurus, C.snowii, he subsequently recognized the elasmosaurid nature of its humerus and coracoids. Thus, he renamed the species E.snowii. A second specimen, discovered by Elias West in 1890, was also assigned by him to E.snowii. In 1943 Welles moved E.snowii to its own genus, Styxosaurus, where the species has remained. However, the West specimen was assigned to Thalassiosaurus ischiadicus (see below) by Welles in 1952; Carpenter returned it to S.snowii in 1999. Williston also reassigned the species E.ischiadicus from the genus Polycotylus, where he had initially placed it when he named it in 1903. The type remains were discovered by him in the same 1874 expedition with Mudge. Williston assigned another specimen discovered by Mudge and H.A. Brous in 1876. In 1943 both specimens were assigned to the new genus Thalassiosaurus by Welles, who then assigned the latter to the new genus and species Alzadasaurus kansasensis in 1952. Glenn Storrs considered both to be indeterminate elasmosaurids in 1999; in the same year, Carpenter assigned both to Styxosaurus snowii.
An elasmosaurid specimen was found by Handel Martin in Logan County, Kansas in 1889. Williston named this as a new species, E.(?)marshii. He bore reservations about its referral to the genus, and he recognized that it possibly pertained to another genus. In 1943 Welles moved E.(?)marshii to a genus of its own, Thalassonomosaurus; however, Carpenter sunk T.marshii into Styxosaurus snowii in 1999. Another species, E.nobilis, was named by Williston from very large remains discovered by Mudge in 1874 in Jewell County, Kansas. Welles named E.nobilis as a species of Thalassonomosaurus, T.nobilis, in 1943, but it too was considered to be part of S.snowii by Carpenter. Finally, two exceptionally large dorsal vertebrae collected by Charles Sternberg in 1895 were named E.sternbergii by Williston, but were considered indeterminate by Storrs. Williston mentioned three additional Elasmosaurus species, which he would figure and describe at a later date. He again made reference to a new species of Elasmosaurus, from Kansas, in 1908.
Several Russian species, based on poorly preserved vertebral remains, were assigned to Elasmosaurus by Nikolay N. Bogolubov in 1911. One was E.helmerseni, which was first described by W.Kiprijanoff in 1882 from Maloje Serdoba, Saratov, as Plesiosaurus helmerseni. Some material from Scania, Sweden, was assigned to P. helmerseni in 1885 by H.Schröder. Vertebral and limb remains from Kursk initially assigned by Kiprijanoff to P.helmerseni were also moved by Bogolubov to the new species E.kurskensis, which he considered to be "identical with Elasmosaurus or related to it". He also named E.orskensis, based on "very large" neck and tail vertebra remains from Konopljanka, Orenburg; and E.serdobensis, based on a single neck vertebra from Maloje Serdoba. However, the validity of all these species has been questioned. Welles considered E.kurskensis as an indeterminate plesiosaur in 1962. Persson noted in a 1959 review of the Swedish "E."helmerseni material that, while the species was probably closely related to Elasmosaurus proper, it was too fragmentary for this hypothesis to be assessed; he later remarked in 1963 that, regarding the latter three species, "their generic and specific definition is questionable", although he declined to specifically label them as invalid on account of not having seen the fossil material. Similarly, in 1999, Evgeniy Pervushov, Maxim Arkhangelsky, and Alexander Ivanov considered E.helmerseni to be an indeterminate elasmosaurid. In 2000 Storrs, Archangelsky, and Vladimir Efimov concurred with Welles on E.kurskensis, and labelled E.orskensis and E.serdobensis as indeterminate elasmosaurids.
Two additional Russian species were described by subsequent authors. Anatoly Riabinin described a single phalanx from a flipper in 1915 as E.(?)sachalinensis; the species was named after the island of Sakhalin, where N.N. Tikhonovich found it in 1909. However, this specimen cannot be identified more specifically than an indeterminate elasmosaurid, which was followed by Persson and Pervushov and colleagues. Storrs, Arkhangelsky, and Efimov were less specific, labelling it as an indeterminate plesiosaur; this classification was followed by Alexander Averianov and Vasilii Popov in 2005. Then, in 1916, Pavel A. Pravoslavlev named E.amalitskii from the Don River region, based on a specimen containing vertebrae, limb girdles, and limb bones. Persson considered it a valid species, and a relatively large member of the elasmosaurids; however, like E.(?)sachalinensis, Pervushov and colleagues considered E.amalitskii an indeterminate elasmosaurid.
In a 1918 review of the geographic distribution and evolution of Elasmosaurus, Pravoslavlev provisionally assigned three other previously named species to Elasmosaurus; his taxonomic opinions have not been widely followed. One of these was E.chilensis, based on the Chilean Plesiosaurus chilensis named from a single tail vertebra by Claude Gay in 1848. In a work published in 1889, Richard Lydekker assigned this species to Cimoliasaurus. Wilhelm Deecke moved chilensis to Pliosaurus in 1895, a classification which was acknowledged by Pravoslavlev. Edwin Colbert later assigned the type vertebra in 1949 to a pliosauroid, and also assigned other assigned remains to indeterminate elasmosauroids; the type vertebra was recognized as potentially belonging to Aristonectes parvidens by José O'Gorman and colleagues in 2013. Another was E.haasti, originally Mauisaurus haasti, named by James Hector in 1874 based on remains found in New Zealand. Although its validity was supported for a considerable time, M.haasti is regarded as a nomen dubium as of 2017. Pravoslavlev recognized another species from New Zealand, E.hoodii, named by Owen in 1870 as Plesiosaurus hoodii based on a neck vertebra. Welles recognized it as a nomen dubium in 1962; Joan Wiffen and William Moisley concurred in a 1986 review of New Zealand plesiosaurs.
In 1949 Welles named a new species of Elasmosaurus, E.morgani. It was named from a well-preserved skeleton found in Dallas County, Texas. However, part of the specimen was accidentally thrown out during the relocation of the Southern Methodist University's paleontological collections. Welles recognized E.morganis similarity to E.platyurus in its shoulder girdle, but maintained it as a separate species due to its shorter neck and more robust rear neck vertebrae. In 1997 Carpenter reconsidered the differences between the two species, and found them sufficient to place E.morgani in its own genus, which he named Libonectes. Despite its reassignment and the loss of its material, L.morgani is often considered an archetypal elasmosaurid. Data based on these lost elements were unquestionably accepted in subsequent phylogenetic analyses, until a redescription of the surviving elements was published by Sachs and Benjamin Kear in 2015.
Persson assigned another species to Elasmosaurus alongside his 1959 description of "E."helmerseni remains from Sweden, namely E.(?)gigas. It was based on Schröder's Pliosaurus(?) gigas, named in 1885 from two dorsals; one was found in Prussia, the other in Scania. While they were incomplete, Persson recognized that their proportions and the shape of their articular ends differed greatly from pliosauroids, and instead agreed well with elasmosaurids. Given that, at the time of Persson's writing, "there [was] nothing to contradict that they are nearest akin to Elasmosaurus", he assigned them to Elasmosaurus "with hesitation". Theodor Wagner had previously assigned gigas to Plesiosaurus in 1914. As of 2013, this questionable attribution remains unchanged. Another species from Russia, E.antiquus, was named by Dubeikovskii and Ochev in 1967 from the Kamsko-Vyatsky phosphorite quarry, but Pervushov and colleagues in 1999, followed by Storrs and colleagues in 2000, reinterpreted it as an indeterminate elasmosaurid.
Classification
Though Cope had originally recognized Elasmosaurus as a plesiosaur, in an 1869 paper he placed it, with Cimoliasaurus and Crymocetus, in a new order of sauropterygian reptiles. He named the group Streptosauria, or "reversed lizards", due to the orientation of their individual vertebrae supposedly being reversed compared to what is seen in other vertebrate animals. He subsequently abandoned this idea in his 1869 description of Elasmosaurus, where he stated he had based it on Leidy's erroneous interpretation of Cimoliasaurus. In this paper, he also named the new family Elasmosauridae, containing Elasmosaurus and Cimoliasaurus, without comment. Within this family, he considered the former to be distinguished by a longer neck with compressed vertebrae, and the latter by a shorter neck with square, depressed vertebrae.
In subsequent years, Elasmosauridae came to be one of three groups in which plesiosaurs were classified, the others being the Pliosauridae and Plesiosauridae (sometimes merged into one group). Charles Andrews elaborated on differences between elasmosaurids and pliosaurids in 1910 and 1913. He characterized elasmosaurids by their long necks and small heads, as well as by their rigid and well-developed scapulae (but atrophied or absent clavicles and interclavicles) for forelimb-driven locomotion. Meanwhile, pliosaurids had short necks but large heads, and used hindlimb-driven locomotion. Although the placement of Elasmosaurus in the Elasmosauridae remained uncontroversial, opinions on the relationships of the family became variable over subsequent decades. Williston created a revised taxonomy of plesiosaurs in 1925.
In 1940 Theodore White published a hypothesis on the interrelationships between different plesiosaurian families. He considered Elasmosauridae to be closest to the Pliosauridae, noting their relatively narrow coracoids as well as their lack of interclavicles or clavicles. His diagnosis of the Elasmosauridae also noted the moderate length of the skull (i.e., a mesocephalic skull); the neck ribs having one or two heads; the scapula and coracoid contacting at the midline; the blunted rear outer angle of the coracoid; and the pair of openings (fenestrae) in the scapula–coracoid complex being separated by a narrower bar of bone compared to pliosaurids. The cited variability in the number of heads on the neck ribs arises from his inclusion of Simolestes to the Elasmosauridae, since the characteristics of "both the skull and shoulder girdle compare more favorably with Elasmosaurus than with Pliosaurus or Peloneustes." He considered Simolestes a possible ancestor of Elasmosaurus. Oskar Kuhn adopted a similar classification in 1961.
Welles took issue with White's classification in his 1943 revision of plesiosaurs, noting that White's characteristics are influenced by both preservation and ontogeny. He divided plesiosaurs into two superfamilies, the Plesiosauroidea and Pliosauroidea, based on neck length, head size, ischium length, and the slenderness of the humerus and femur (the propodialia). Each superfamily was further subdivided by the number of heads on the ribs, and the proportions of the epipodialia. Thus, elasmosaurids had long necks, small heads, short ischia, stocky propodialia, single-headed ribs, and short epipodialia. Pierre deSaint-Seine in 1955 and Alfred Romer in 1956 both adopted Welles' classification. In 1962 Welles further subdivided elasmosaurids based on whether they possessed pelvic bars formed from the fusion of the ischia, with Elasmosaurus and Brancasaurus being united in the subfamily Elasmosaurinae by their sharing of completely closed pelvic bars.
Carpenter's 1997 phylogenetic analysis of plesiosaurs challenged the traditional subdivision of plesiosaurs based on neck length. While polycotylids had previously been part of the Pliosauroidea, Carpenter moved polycotylids to become the sister group of the elasmosaurids based on similarities, thus implying that polycotylids and pliosauroids evolved their short necks independently. The content of Elasmosauridae also received greater scrutiny. Since its initial assignment to the Elasmosauridae, the relationships of Brancasaurus had been considered well supported, and an elasmosaurid position was recovered by O'Keefe's 2004 analysis and Franziska Großmann's 2007 analysis. However, Ketchum and Benson's analysis instead included it in the Leptocleidia, and its inclusion in that group has remained consistent in subsequent analyses. Their analysis also moved Muraenosaurus to the Cryptoclididae, and Microcleidus and Occitanosaurus to the Plesiosauridae; Benson and Druckenmiller isolated the latter two in the group Microcleididae in 2014, and considered Occitanosaurus a species of Microcleidus. These genera had all previously been considered to be elasmosaurids by Carpenter, Großmann, and other researchers.
Within the Elasmosauridae, Elasmosaurus itself has been considered a "wildcard taxon" with highly variable relationships. Carpenter's 1999 analysis suggested that Elasmosaurus was more basal (i.e. less specialized) than other elasmosaurids with the exception of Libonectes. In 2005 Sachs suggested that Elasmosaurus was closely related to Styxosaurus, and in 2008 Druckenmiller and Russell placed it as part of a polytomy with two groups, one containing Libonectes and Terminonatator, the other containing Callawayasaurus and Hydrotherosaurus. Ketchum and Benson's 2010 analysis included Elasmosaurus in the former group. Benson and Druckenmiller's 2013 analysis (below, left) further removed Terminonatator from this group and placed it as one step more derived. In Rodrigo Otero's 2016 analysis based on a modification of the same dataset (below, right), Elamosaurus was the closest relative of Albertonectes, forming the Styxosaurinae with Styxosaurus and Terminonatator. Danielle Serratos, Druckenmiller, and Benson could not resolve the position of Elasmosaurus in 2017, but they noted that Styxosaurinae would be a synonym of Elasmosaurinae if Elasmosaurus did fall within the group. In 2020, O'Gorman formally synonymized Styxosaurinae with Elasmosaurinae based on the inclusion of Elasmosaurus within the group, and also provided a list of diagnostic characteristics for the clade.
Topology A: Benson et al. (2013)
Topology B: Otero (2016), with clade names following O'Gorman (2020)
Paleobiology
Elasmosaurids were fully adapted to life in the ocean, with streamlined bodies and long paddles that indicate they were active swimmers. The unusual body structure of elasmosaurids would have limited the speed at which they could swim, and their paddles may have moved in a manner similar to the movement of oars rowing, and due to this, could not twist and were thus held rigidly. Plesiosaurs were even believed to have been able to maintain a constant and high body temperature (homeothermy), allowing for sustained swimming.
A 2015 study concluded that locomotion was mostly done by the fore-flippers while the hind-flippers functioned in maneuverability and stability; a 2017 study concluded that the hind-flippers of plesiosaurs produced 60% more thrust and had 40% more efficiency when moving in harmony with the fore-flippers. The paddles of plesiosaurs were so rigid and specialized for swimming that they could not have come on land to lay eggs like sea turtles. Therefore, they probably gave live-birth (viviparity) to their young like some species of sea snakes. Evidence for live-birth in plesiosaurs is provided by the fossil of an adult Polycotylus with a single fetus inside.
Elasmosaurid remains provide some evidence they were preyed upon. A humerus of an unidentified subadult elasmosaurid was found with bite marks matching the teeth of the shark Cretoxyrhina, while a crushed Woolungasaurus skull has tooth-marks matched to the pliosaur Kronosaurus.
Neck movement and function
Cope, in 1869, compared the build and habits of Elasmosaurus with those of a snake. Although he suggested that the vertebral column of the trunk did not allow for much vertical movement due to the elongated neural spines which nearly form a continuous line with little space between adjacent vertebrae, he envisaged the neck and tail to have been much more flexible: "The snake-like head was raised high in the air, or depressed at the will of the animal, now arched swan-like preparatory to a plunge after a fish, now stretched in repose on the water or deflexed in exploring the depths below".
Although followed by many common media depictions, more recent research showed that Elasmosaurus was incapable of raising anything more than its head above the water. The weight of its long neck placed the center of gravity behind the front flippers. Thus, Elasmosaurus could have raised its head and neck above the water only when in shallow water, where it could rest its body on the bottom. Also, the weight of the neck, the limited musculature, and the limited movement between the vertebrae would have prevented Elasmosaurus from raising its head and neck very high. The head and shoulders of the Elasmosaurus probably acted as a rudder. If the animal moved the anterior part of the body in a certain direction, it would cause the rest of the body to move in that direction. Thus, Elasmosaurus would have been unable to swim in one direction while moving its head and neck either horizontally or vertically in a different direction.
One study found that the necks of elasmosaurids were capable of 75–177˚ of ventral movement, 87–155° of dorsal movement, and 94–176° of lateral movement, depending on the amount of tissue between the vertebrae, which probably increased in rigidness towards the back of the neck. The researchers concluded that lateral and vertical arches and shallow S-shaped curves were feasible in contrast to the "swan-like" S-shape neck postures that required more than 360° of vertical flexion.
The exact function of the neck of elasmosaurids is unknown, though it may have been important for hunting. It has also been suggested that the long necks of plesiosaurs served as a snorkel and allowed them to breathe air while the body remained underwater. This is disputed as there would be large hydrostatic pressure differences, particularly for the extremely long-necked elasmosaurids. The neck anatomy of elasmosaurids was capable of making a gentle slope to allow them to breathe at the surface but would have required them to engage in energy-expensive swimming at the sub-surface. In addition, the longer neck would also have increased dead space, and the animals may have required larger lungs. The neck could have had other vulnerabilities, for example being a target for predators.
Simulation of water flow on 3D models showed that more elongated necks, such as those of elasmosaurids, did not increase drag force while swimming compared to shorter necked plesiosaurs. On the other hand, bending the neck sideways did increase drag force, more so in forms with very long necks. Another study found the long necks of elasmosaurs would normally increase drag during forward swimming but this was cancelled out by their large torsos, and hence large body sizes may have facilitated the evolution of longer necks.
Feeding
In 1869 Cope noted that scales and teeth of six species of fish had been discovered directly beneath the vertebrae of the Elasmosaurus holotype, and theorized that these fish would have had formed the diet of the animal. From these remains, Cope named a new species of barracuda, Sphyraena carinata.
The flexion ranges of Elasmosaurus necks would have allowed the animal to employ a number of hunting methods including "benthic grazing", which would have involved swimming close to the bottom and using the head and neck to dig for prey on the sea floor. Elasmosaurids may also have been active hunters in the pelagic zone, retracting their necks to launch a strike or using side-swipe motions to stun or kill prey with their laterally projected teeth (like sawsharks). It has also been suggested that the predatory abilities of elasmosaurids have been underestimated; their large skulls, big jaw-muscles, strong jaws, and long teeth indicate they could prey on animals between and long, as indicacted by stomach contents including those of sharks, fish, mosasaurs, and cephalopods.
It is possible that Elasmosaurus and its kin stalked schools of fish, concealing themselves below and moving the head slowly up as they approached. The eyes of the animal were at the top of the head and allowed them to see directly upward. This stereoscopic vision would have helped it to find small prey. Hunting from below would also have been possible, with prey silhouetted in the sunlight while concealed in the dark waters below. Elasmosaurids probably ate small bony fish and marine invertebrates, as their small, non-kinetic skulls would have limited the size of the prey they could eat. Also, with their long, slender teeth adapted for seizing prey and not tearing, elasmosaurids most certainly swallowed their prey whole.
Although elasmosaurids are commonly found with several gastroliths, Elamosaurus has only been found uncontroversially with a pebble lodged in the neural arch of one of its hindmost tail-vertebrae. A specimen of the closely related Styxosaurus contained fragmented fish bones and stones in the abdominal region behind the pectoral girdle. The fish remains were identified as Enchodus and other clupeomorph fish. The stones match rock from away from where the specimen was found. Several different functions have been proposed for gastroliths, including aiding in digestion, mixing food content, mineral supplementation, and storage and buoyancy control.
Paleoecology
Elasmosaurus is known from the Sharon Springs Member of the Campanian-age Upper Cretaceous Pierre Shale formation of western Kansas, which dates to about 80.5million years ago. The Pierre Shale represents a period of marine deposition from the Western Interior Seaway, a shallow continental sea that submerged much of central North America during the Cretaceous. At its largest, the Western Interior Seaway stretched from the Rockies east to the Appalachians, some wide. At its deepest, it may have been only deep. Two great continental watersheds drained into it from east and west, diluting its waters and bringing resources in eroded silt that formed shifting river delta systems along its low-lying coasts. There was little sedimentation on the eastern margin of the Seaway; the western margin accumulated a thick pile of sediments eroded from the western land mass. The western shore was thus highly variable, depending on variations in sea level and sediment supply.
The soft, muddy sea floor probably received very little sunlight, but it teemed with life due to steady rains of organic debris from plankton and other organisms farther up the water column. The bottom was dominated by large Inoceramus clams, which were covered with oysters; there was little biodiversity. Clam shells would have accumulated over the centuries in layers under the sea floor's surface, and would have provided shelter for small fish. Other invertebrates known to have lived in this sea include various species of rudists, crinoids and cephalopods (including squids and ammonites).
Large fish known to have inhabited the sea include the bony fishes Pachyrhizodus, Enchodus, Cimolichthys, Saurocephalus, Saurodon, Gillicus, Ichthyodectes, Xiphactinus, Protosphyraena and Martinichthys; and the sharks Cretoxyrhina, Cretolamna, Scapanorhynchus, Pseudocorax and Squalicorax. In addition to Elasmosaurus, other marine reptiles present include fellow plesiosaurs Libonectes, Styxosaurus, Thalassomedon, Terminonatator, Polycotylus, Brachauchenius, Dolichorhynchops and Trinacromerum; the mosasaurs Mosasaurus, Halisaurus, Prognathodon, Tylosaurus, Ectenosaurus, Globidens, Clidastes, Platecarpus and Plioplatecarpus; and the sea turtles Archelon, Protostega, Porthochelys and Toxochelys. The flightless aquatic bird Hesperornis also made its home there. The pterosaurs Pteranodon and Nyctosaurus, and the bird Ichthyornis, are also known far from land.
| Biology and health sciences | Prehistoric marine reptiles | Animals |
1264239 | https://en.wikipedia.org/wiki/Reaction%20coordinate | Reaction coordinate | In chemistry, a reaction coordinate is an abstract one-dimensional coordinate chosen to represent progress along a reaction pathway. Where possible it is usually a geometric parameter that changes during the conversion of one or more molecular entities, such as bond length or bond angle. For example, in the homolytic dissociation of molecular hydrogen, an apt choice would be the coordinate corresponding to the bond length. Non-geometric parameters such as bond order are also used, but such direct representation of the reaction process can be difficult, especially for more complex reactions.
In computer simulations collective variables are employed for a target-oriented sampling approach. Plain simulations fail to capture so called rare events, because they are not feasible to occur in realistic computation times. This often stems from to high energy barriers separating the reactants from products, or any two states of interest. A collective variable is as the name states only a set, a collection, of individual variables () contracted into one:
,
with a transformation matrix. The collective variables reduce many variables to a lower-dimensional set of variables, that still describe the crucial characteristics of the system. Many collective variables than span the reaction coordinate with a continuous function :
.
An example is the complexation of two molecules. The distance between both of them is the collective variable, where the atomic positions are the individual variables and the reaction coordinate would be the full path of association and dissociation. By applying a bias to the collective variables the simulation can be 'steered' towards the desired destination. These kinds of simulations are called enhanced simulations.
Special collective variables that help to distinguish reactants from products are also known as order parameters, terminology that originates in work on phase transitions. Reaction coordinates are special order parameters that describe the entire pathway from reactants through transition states and on to products. Depending on the application, reaction coordinates may be defined by using chemically intuitive variables like bond lengths, or splitting probabilities (also called committors), or using the eigenfunction corresponding to the reactant-to-product transition as a progress coordinate.
A reaction coordinate parameterizes reaction process at the level of the molecular entities involved. It differs from extent of reaction, which measures reaction progress in terms of the composition of the reaction system.
(Free) energy is often plotted against reaction coordinate(s) to demonstrate in schematic form the potential energy profile (an intersection of a potential energy surface) associated with the reaction.
In the formalism of transition-state theory the reaction coordinate for each reaction step is one of a set of curvilinear coordinates obtained from the conventional coordinates for the reactants, and leads smoothly among configurations, from reactants to products via the transition state. It is typically chosen to follow the path defined by potential energy gradient – shallowest ascent/steepest descent – from reactants to products.
| Physical sciences | Kinetics | Chemistry |
1264353 | https://en.wikipedia.org/wiki/Tantalite | Tantalite | The mineral group tantalite [(Fe, Mn)Ta2O6] is the primary source of the chemical element tantalum, a corrosion (heat and acid) resistant metal. It is chemically similar to columbite, and the two are often grouped together as a semi-singular mineral called coltan or "columbite-tantalite" in many mineral guides. However, tantalite has a much greater specific gravity than columbite (8.0+ compared to columbite's 5.2). Iron-rich tantalite is the mineral tantalite-(Fe) or ferrotantalite and manganese-rich is tantalite-(Mn) or manganotantalite.
Tantalite is also very close to tapiolite. Those minerals have the same chemical composition, but different crystal symmetry: orthorhombic for tantalite and tetragonal for tapiolite.
Tantalite is black to brown in both color and streak. Manganese-rich tantalites can be brown and translucent.
Occurrence
Tantalite occurs in granitic pegmatites that are rich in rare-elements, and in placer deposits derived from such rocks. It has been found in Australia, Brazil, Canada, Colombia (Guainía and Vichada), Egypt, northern Europe, Madagascar, Namibia, Nigeria, Rwanda, The Democratic Republic of Congo, the United States (California, Colorado, Maine, and Virginia), and Zimbabwe. Brazil has the world's largest reserve of tantalite (52.1%).
Applications
The tantalum metal extracted from tantalite is used in alloys for strength and higher melting points, in glass to increase the index of refraction, and in surgical steel, as it is non-reactive and non-irritating to body tissues. Much like glass, it is not suitable for use in hydrofluoric acid and strong hot alkali applications.
Sustainability
The mining of tantalite causes many environmental and social problems in the Democratic Republic of Congo.
| Physical sciences | Minerals | Earth science |
1264869 | https://en.wikipedia.org/wiki/Albendazole | Albendazole | Albendazole is a broad-spectrum antihelmintic and antiprotozoal agent of the benzimidazole type. It is used for the treatment of a variety of intestinal parasite infections, including ascariasis, pinworm infection, hookworm infection, trichuriasis, strongyloidiasis, taeniasis, clonorchiasis, opisthorchiasis, cutaneous larva migrans, giardiasis, and gnathostomiasis, among other diseases.
Common side effects include nausea, abdominal pain, and headache. Rare but potentially serious side effects include bone marrow suppression which usually improves on discontinuing the medication. Liver inflammation has been reported and those with prior liver problems are at greater risk. It is pregnancy category D in Australia, meaning it may cause harm if taken by pregnant women.
Albendazole was developed in 1975. It is on the World Health Organization's List of Essential Medicines.
Medical uses
Albendazole is an effective treatment for:
Flatworms
Clonorchiasis
Fasciolosis
Opisthorchiasis
Cestodes (tapeworms), as an alternative to praziquantel or niclosamide for adult beef tapeworms and as an alternative to praziquantel for pork tapeworms. It is also given for infections by T. crassiceps. Though praziquantel is often better at treating tapeworm infections, albendazole is used more often in endemic countries due to being cheaper and having a broader spectrum.
Cysticercosis (especially neurocysticercosis), which is caused by the larval form of the pork tapeworm (i.e. albendazole is the drug of choice for larval pork tapeworms, but not adult pork tapeworms). Old cysts are not affected.
Echinococcosis of the liver, lung, and peritoneum (caused by the larval form of the dog tapeworm, or of the alveoli (caused by E. multilocularis) when surgical excision is not possible. Alveolar and cystic echinococcosis may require lifelong treatment with albendazole, which only prevents the parasites from growing and reproducing rather than killing them.
Nematodes
Anatrichosomiasis
Angiostrongyliasis
Anisakiasis
Ascariasis, which can be cured with a single dose of albendazole.
Baylisascariasis, caused by the raccoon roundworm. Albendazole can achieve good results (95–100% efficacy after a 10-day course of treatment) if treatment is initiated within 72 hours of ingestion of the egg-containing raccoon feces. Corticosteroids are sometimes added in cases of eye and CNS infections.
Pinworm infection
Filariasis; since albendazole's disintegration of the microfilariae ("pre-larva") can cause an allergic reaction, antihistamines or corticosteroids are sometimes added to treatment. In cases of lymphatic filariasis (elephantiasis) caused by Wuchereria bancrofti or Brugia malayi, albendazole is sometimes given as an adjunct to ivermectin or diethylcarbamazine in order to suppress microfilaremia. It can also be given for Loa loa filariasis as an adjunct or replacement to diethylcarbamazine. Albendazole has an embryotoxic effect on Loa loa adults and thus slowly reduces microfilaremia.
Gnathostomiasis when caused by Gnathostoma spinigerum. Albendazole has a similar effectiveness to ivermectin in these cases, though it needs to be given for 21 days rather than the 2 days needed for ivermectin.
Gongylonemiasis
Hepatic capillariasis caused by Capillaria hepatica
Hookworm infections, including cutaneous larva migrans caused by hookworms of genus Ancylostoma. A single dose of albendazole is sufficient to treat intestinal infestations by A. duodenale or Necator americanus.
Intestinal capillariasis, as an alternative to mebendazole
Mansonelliasis when caused by Mansonella perstans. Albendazole is effective against adult worms but not against the immature microfilariae.
Oesophagostomumiasis, when caused by Oesophagostomum bifurcum
Strongyloidiasis, as an alternative to ivermectin or thiabendazole. Albendazole can be given with diethylcarbamazine to lower microfilaremia levels.
Toxocariasis, also called "visceral larva migrans", when caused by the dog roundworm Toxocara canis or cat roundworm T. catis. Corticosteroids can be added in severe cases, and surgery might be required to repair secondary damage.
Trichinosis, when caused by Trichinella spiralis or T. pseudospiralis. Albendazole has a similar efficacy to thiabendazole, but fewer side effects. It works best when given early, acting on the adult worms in the intestine before they generate larva that can penetrate the muscle and cause a more widespread infection. Corticosteroids are sometimes added on to prevent inflammation caused by dying larva.
Trichostrongyliasis, as an alternative to pyrantel pamoate. A single dose is sufficient for treatment.
Trichuriasis (whipworm infection), sometimes considered as an alternative to mebendazole and sometimes considered to be the drug of choice. Only a single dose of albendazole is needed. It can also be given with ivermectin.
Giardiasis, as an alternative or adjunct to metronidazole, especially in children
Microsporidiosis, including ocular microsporidiosis caused by Encephalitozoon hellem or E. cuniculi, when combined with topical fumagillin
Granulomatous amoebic encephalitis, when caused by the amoeba Balamuthia mandrillaris, in combination with miltefosine and fluconazole
Arthropods
Crusted scabies, when combined with topical crotamiton and salicylic acid
Head lice infestation, though ivermectin is much better
Intestinal myiasis
Though albendazole is effective in treating many diseases, it is only FDA-approved for treating hydatid disease caused by dog tapeworm larvae and neurocysticercosis caused by pork tapeworm larvae.
Pregnancy
Albendazole is a pregnancy class D drug in Australia. It is contraindicated in the first trimester of pregnancy, and should be avoided up to one month before conception. While studies in pregnant rats and rabbits have shown albendazole to be teratogenic, albendazole has been found to be safe in humans during the second and third trimesters. It can, however, possibly cause infantile eczema when given during pregnancy.
In pregnant dogs, albendazole use has led to puppies with reduced weight and with cleft palates. Birds have lower rates of laying eggs and hatching when given albendazole.
Albendazole sulfoxide is secreted into breast milk at around 1.5% of the maternal dose, though oral absorption is poor enough that it is unlikely to affect nursing infants.
Contraindications
Hypersensitivity to the benzimidazole class of compounds contraindicates its use.
Side effects
The most common side effects of albendazole are experienced by over 10% of people and include headache and abnormal liver function. Elevation of liver enzymes occurs in 16% of patients receiving treatment specifically for hydatid disease and goes away when treatment ends. Liver enzymes usually increase to two to four times the normal levels (a mild to moderate increase). An estimated 1–10% of people experience abdominal pain, nausea or vomiting, dizziness or vertigo, increased intracranial pressure, meningeal signs, temporary hair loss, and fever. The headache, nausea, and vomiting are thought to be caused by the sudden destruction of cysticerci (tapeworm larvae), which causes acute inflammation. Fewer than 1% of people get hypersensitivity reactions such as rashes and hives, leukopenias (drop in white blood cell levels) such as agranulocytosis and granulocytopenia, thrombocytopenia (reduced platelet count), pancytopenia (drop in white blood cells, red blood cells, and platelets), hepatitis, acute liver failure, acute kidney injury, irreversible bone marrow suppression, and aplastic anemia.
Side effects can be different when treating for hydatid disease versus neurocysticercosis: for example, those being treated for the former are more likely to experience elevated liver enzymes and abdominal pain, while those being treated for the latter are more likely to experience headache. Treating hydatid disease can also unmask undiagnosed neurocysticercosis. People receiving albendazole for the treatment of neurocysticercosis can have neurological side effects such as seizures, increased intracranial pressure, and focal signs caused by the inflammatory reaction that occurs when parasites in the brain are killed. Steroids and anticonvulsants are often given with albendazole when treating neurocysticercosis to avoid these effects. Those being treated for retinal neurocysticercosis can face retinal damage if they are not first checked for ocular cysticeri, since changes to existing lesions in the eye by albendazole can cause permanent blindness.
Overdose
Because of its low solubility, albendazole often cannot be absorbed in high enough quantities to be toxic. The oral LD50 of albendazole in rats was found to be 2,500 mg/kg. It takes 20 times the normal dose to kill a sheep, and 30 times the normal dose to kill cattle. Overdose affects the liver, testicles, and gastrointestinal tract (GI tract) the most. It can manifest with lethargy, loss of appetite, vomiting, diarrhea, intestinal cramps, dizziness, convulsions, and sleepiness. There is no specified antidote.
Interactions
The antiepileptics carbamazepine, phenytoin, and phenobarbital lower the plasma concentration and half-life of albendazole sulfoxide's R(+) enantiomer.
The antacid cimetidine heightens serum albendazole concentrations, increases the half-life of albendazole, and doubles albendazole sulfoxide levels in bile. It was originally thought to work by increasing albendazole bioavailability directly; however, it is now known that cimetidine inhibits the breakdown of albendazole sulfoxide by interfering with CYP3A4. The half-life of albendazole sulfoxide thus increases from 7.4 hours to 19 hours. This might be a helpful interaction on more severe cases, because it boosts the potency of albendazole. Paradoxically, cimetidine also inhibits the absorption of albendazole by reducing gastric acidity.
Several other interactions exist. Corticosteroids increase the steady-state plasma concentration of albendazole sulfoxide; dexamethasone, for example, can increase the concentration by 56% by inhibiting the elimination of albendazole sulfoxide. The anti-parasitic praziquantel increases the maximum plasma concentration of albendazole sulfoxide by 50%, and the anti-parasitic levamisole increases the AUC (total drug exposure) by 75%. Grapefruit inhibits the metabolism of albendazole within the intestinal mucosa. Finally, long-term administration of the antiretroviral ritonavir, which works as a CYP3A4 inhibitor, decreases the maximum concentration of albendazole in the plasma as well as the AUC.
Pharmacology
Mechanism of action
As a vermicide, albendazole causes degenerative alterations in the intestinal cells of the worm by binding to the colchicine-sensitive site of β-tubulin, thus inhibiting its polymerization or assembly into microtubules (it binds much better to the β-tubulin of parasites than that of mammals). Albendazole leads to impaired uptake of glucose by the larval and adult stages of the susceptible parasites, and depletes their glycogen stores. Albendazole also prevents the formation of spindle fibers needed for cell division, which in turn blocks egg production and development; existing eggs are prevented from hatching. Cell motility, maintenance of cell shape, and intracellular transport are also disrupted. At higher concentrations, it disrupts the helminths' metabolic pathways by inhibiting metabolic enzymes such as malate dehydrogenase and fumarate reductase, with inhibition of the latter leading to less energy produced by the Krebs cycle. Due to diminished ATP production, the parasite is immobilized and eventually dies.
Some parasites have evolved some resistance to albendazole by having a different set of amino acids constitute β-tubulin, decreasing the binding affinity of albendazole. Some parasites (especially filarial nematodes) live in symbiosis with Wolbachia, a type of intracellular parasite bacteria. In such cases the Wolbachia are necessary to the survival of the parasitic worms. Elimination of Wolbachia from these filarial nematodes generally results in either death or sterility of the host nematode.
Pharmacokinetics
To target intestinal parasites, which is the most common indication for prescription, albendazole is taken on an empty stomach to stay within the gut.
Oral absorption of albendazole varies among species, with 1–5% of the drug being successfully absorbed in humans, 20–30% in rats, and 50% in cattle.
The absorption also largely depends on gastric pH. People have varying levels of gastric pHs on empty stomachs, and thus absorption from one person to another can vary wildly when taken without food. Generally, the absorption in the GI tract is poor due to albendazole's low solubility in water. It is, however, better absorbed than other benzimidazole carbamates. Food stimulates gastric acid secretion, lowering the pH and making albendazole more soluble and thus more easily absorbed. Oral absorption is especially increased with a fatty meal, as albendazole dissolves better in lipids, allowing it to cross the lipid barrier created by the mucus surface of the GI tract.
Absorption is also affected by how much of the albendazole is degraded within the small intestine by metabolic enzymes in the villi.
The pharmacokinetics of albendazole differ slightly between men and women: women have a lower oral clearance and volume of distribution, while men have a lower serum peak concentration.
Albendazole undergoes very fast first-pass metabolism in all species, such that the unchanged drug is undetectable in plasma. Most of it is oxidized into albendazole sulfoxide (also known as ricobendazole and albendazole oxide) in the liver by cytochrome P450 oxidases (CYPs) and a flavin-containing monooxygenase (FMO), which was discovered later. In humans, the cytochrome P450 oxidases are thought to include CYP3A4 and CYP1A1, while those in the rats are thought to be CYP2C6 and CYP2A1.
Oxidation to albendazole sulfoxide by FMO produces R(+) enantiomers, while oxidation the cytochromes and by some enzymes in the gut epithelium produce S(−). Different species produce the R(+) and S(−) enantiomers in different quantities; humans, dogs, and most other species produce the R(+) enantiomer more (with the human AUC ratio being 80:20). Compared to the S(−) enantiomer, the R(+) has greater pharmacological activity, lasts longer in the bloodstream, is found in higher concentrations in the infected host tissues, and is found in higher concentrations within the parasites themselves. Some albendazole is also converted to hydroxyalbendazole, mainly by CYP2J2.
For systemic parasites, albendazole acts as a prodrug, while albendazole sulfoxide reaches systemic circulation and acts as the real antihelminthic. Albendazole sulfoxide is able to cross the blood–brain barrier and enter the cerebrospinal fluid at 43% of plasma concentrations; its ability to enter the central nervous system allows it to treat neurocysticercosis.
Albendazole sulfoxide is converted to the inactive albendazole sulfone by cytochrome P450 oxidases, thought to include CYP3A4 or CYP2C. Other inactive metabolites include: 2-aminosulfone, ω-hydroxysulfone, and β-hydroxysulfone. The major final metabolites that are excreted by humans are:
methyl [5-(propylsulfonyl-1H-benzimidazol-2-yl)] carbamate,
methyl [6-hydroxy 5-(n-propylsulfonyl)-1H-benzimidazole-2-yl)] carbamate,
methyl [5-(n-propylsulfinyl)-1H-benzimidazole-2-yl)] carbamate,
5-(n-propylsulfonyl)-1H-benzimidazole-2-yl amine, and
5-(n-propysulfinyl)-1H-benzimidazole-2-yl amine.
There are also some minor hydroxylated sulfated or glucuronidated derivatives. No unchanged albendazole is excreted, as it is metabolized too quickly.
In humans, the metabolites are excreted mostly in bile, with only a small amount being excreted in urine (less than 1%) and feces. In ruminants, 60–70% of the metabolites are excreted in urine.
Like all benzimidazoles, albendazole has no residual effect, and thus protects poorly against reinfestations.
History
Albendazole, patented in 1975, was invented by Robert J. Gyurik and Vassilios J. Theodorides and assigned to SmithKline Corporation. It was introduced in 1977 as an antihelminthic for sheep in Australia, and was registered for human use in 1982.
Society and culture
Economics
The pharmaceutical company Amedra increased the price after purchasing the rights to the drug, instead of lowering it as generics are predicted to do, drawing criticism from patients' rights advocates.
In 2013, GlaxoSmithKline donated 763 million albendazole tablets for the treatment and prevention of parasitic infections in developing countries, bringing the total to over 4 billion tablets donated since 1998.
Brand names
Brand names include: Albenza, Alworm, Andazol, Eskazole, Noworm, Zentel, Alben-G, ABZ, Cidazole, Wormnil etc.
Research
Albendazole and related compounds or metabolites like albendazole sulfone (ALB-SO2) exhibit antibacterial effects via an unknown, possibly FtsZ-related, mechanism. It inhibits division of Wolbachia and Mycobacterium tuberculosis, turning them into a long "filament" shape as they grow and fail to divide. Since Brugia malayi relies on symbiotic Wolbachia, this would mean that albendazole is targeting both the worm and its essential symbioant.
Veterinary use
Albendazole is mainly used in cattle and sheep, but has found some use in cats and dogs as well; it is also used in ratite birds for flagellate parasites and tapeworms. It is also used off-label to treat endoparasites in goats and pigs.
Albendazole has been used as an antihelminthic and for control of flukes in a variety of animal species, including cattle, sheep, goats, swine, camels, dogs, cats, elephants, poultry, and others. Side effects include anorexia in dogs and lethargy, depression, and anorexia in cats, with more than 10% of dogs and cats having anorexia. Of dogs and cats, 1–10% experience elevated liver enzymes, nausea, vomiting, and diarrhea. Less than 1% experience neutropenia or aplastic anemia, though these require a use of at least 5 days. While it is also associated with bone marrow suppression and toxicity in cats and dogs at high doses, albendazole has a higher margin of safety in other species. Thus, it is usually only prescribed in cats and dogs when an infection is present that is resistant to the commonly prescribed metronidazole and fenbendazole.
It is extensively used for ruminant livestock in Latin America. It is marketed for this purpose by Zoetis (formerly Pfizer Animal Health) in numerous countries (including the United States and Canada) as Valbazen in oral suspension and paste formulations; by Interchemie in the Netherlands and elsewhere as Albenol-100; by Channelle Animal Health Ltd. in the United Kingdom as Albex; and by Ravensdown in New Zealand (as Albendazole). Although most formulations are administered orally, Ricomax (ricobendazole, or albendazole sulfoxide) is administered by subcutaneous injection.
Albendazole has greater bioavailability in ruminants: some albendazole sulfoxide, when released back into the rumen, is reduced to albendazole by the resident microbiota, with a preference of the (+) enantiomer being the substrate. Cats and dogs, having no rumen reservoir, sometimes need higher or more frequent doses as compared to ruminants. In dogs, albendazole sulfoxide is detectable in the plasma for less than 12 hours, but in sheep and goats, it remains at measurable levels for around three days.
Meat
The limitations in early pregnancy are due to a limited period during which teratogenic effects may occur. Summarized research data relating to the durations of these preslaughter and early pregnancy periods when albendazole should not be administered are found in US FDA NADA 110-048 (cattle) and 140-934 (sheep). Some data and inferences regarding goats are found in US FDA Supplemental NADA 110-048 (approved January 24, 2008).
Maximum residue limits (MRLs) for albendazole in food, adopted by the FAO/WHO Codex Alimentarius in 1993, are 5000, 5000, 100, and 100 micrograms per kilogram of body weight (μg/kg) for kidney, liver, fat, and muscle, respectively, and 100 μg/L for milk. For analysis purposes, MRLs of various nations may pertain to concentration of a marker substance which has been correlated with concentrations of the administered substance and its metabolized products. For example, in Canada, the marker substance specified by Health Canada is albendazole-2-aminosulfone, for which the MRL in liver of cattle is 200 μg/kg.
There is a 27-day cattle withdrawal time for meat.
| Biology and health sciences | Antiparasitic | Health |
1265218 | https://en.wikipedia.org/wiki/Podzol | Podzol | In soil science, podzols, also known as podosols, spodosols, or espodossolos, are the typical soils of coniferous or boreal forests and also the typical soils of eucalypt forests and heathlands in southern Australia. In Western Europe, podzols develop on heathland, which is often a construct of human interference through grazing and burning. In some British moorlands with podzolic soils, cambisols are preserved under Bronze Age barrows.
Term
Podzol means "under-ash" and is derived from the Russian () + (); the full form is (), meaning "under-ashed soil". The term was first given in mid-1875 by Vasily Dokuchaev. It refers to the common experience of Russian peasants of plowing up an apparent under-layer of ash (leached or E horizon) during first plowing of a virgin soil of that type.
Characteristics
Podzols can occur on almost any parent material but generally derive from either quartz-rich sands and sandstone or sedimentary debris from magmatic rocks, provided there is high precipitation. Most Podzols are poor soils for agriculture due to the sandy portion, resulting in a low level of moisture and nutrients. Some are sandy and excessively drained. Others have shallow rooting zones and poor drainage due to subsoil cementation. A low pH further compounds issues, along with phosphate deficiencies and aluminum toxicity. The best agricultural use of Podzols is for grazing, although well-drained loamy types can be very productive for crops if lime and fertilizer are used.
The E horizon (or Ae in Canadian soil classification system), which is usually thick, is low in Fe and Al oxides and humus. It is formed under moist, cool and acidic conditions, especially where the parent material, such as granite or sandstone, is rich in quartz. It is found under a layer of organic material in the process of decomposition, which is usually thick. In the middle, there is often a thin horizon of . The bleached soil horizon, which always has a higher value than the horizons above and below it, goes over into a red or red-brown horizon (so-called Podzolic B). The colour is strongest in the upper part, and change at a depth of progressively to the part of the soil that is mainly not affected by processes; that is the parent material. The soil profiles are designated by the letters A (topsoil), E (eluviated soil), B (subsoil) and C (parent material).
In some Podzols, the E horizon is absent—either masked by biological activity or obliterated by disturbance. Podzols with little or no E horizon development are often classified as brown Podzolic soils, also called Umbrisols or Umbrepts.
Geographic distribution
Podzols cover about worldwide and are usually found under sclerophyllous woody vegetation. By extent Podzols are most common in temperate and boreal zones of the Northern Hemisphere but they can also be found in other settings including both temperate rainforests and tropical areas.
In South America Podzols occur beneath Nothofagus betuloides forests in Tierra del Fuego.
Podzolization
Podzolization (or Podsolization) is a complex soil formation process by which dissolved organic matter and ions of iron and aluminium, released through weathering of various minerals, form organo-mineral complexes (chelates) and are moved from the upper parts of the soil profile and deposit in the deeper parts of soil. Through this process, the eluvial horizon becomes bleached and of ash-grey colour. The complexes move with percolating water further down to illuviated horizons which are commonly coloured brown, red or black as they accumulate and consist of cemented sesquioxides and/or organic compounds. The podzolization is a typical soil formation process in Podzols.
Preconditions
Podzolization usually occurs under forest or heath vegetation and is common in cool and humid climates as these climates inhibit the activity of soil microbes in the topsoil. Overall, podzolization happens where the decomposition of organic matter is inhibited and as a result, acidic organic surface (mor) layers build up. Under these typically acidic conditions, nutrient deficiency further hampers the microbial degradation of organic complexing agents. Medium to coarse textured soils with base-poor parent material (usually rich in quartz) also promote podzolization, as they encourage percolating water flow.
Key steps
The soil-forming process of podzolization can be broken down into two main steps:
Mobilization and translocation of organic matter, Fe and Al from the surface horizon, and
Immobilization and stabilization of organic matter, Fe and Al into the subsoil.
In the topsoil of acidic soils, organic matter (mostly from plant litter, the humus layer and root exudates) together with Al- and Fe-ions, form organo-mineral complexes. These soluble chelates then relocate with percolating water from the A (or E horizon) to the B horizon. As a result of this, the E horizon (or Ae horizon in the Canadian system of soil classification) is left bleached and ash-grey in colour, while the B horizon becomes enriched with relocated organo-mineral complexes. The colour of B horizon is consequently red, brown or black, depending on the dominance of metal ions or organic matter. Usually, the boundary between the B and eluvial Ae (or E) horizon is very distinct, and sometimes a hardpan (or Ortstein) can form, as the relocated Fe and Al and organic matter increase mineral particles, cementing them into this compacted layer.
There are several reasons why these organo-mineral complexes immobilize in the B horizon: If during the eluviation process more Al- or Fe-ions bind to the organic compounds, the complex can flocculate as the solubility of it decreases with increasing metal to carbon ratio. Apart from that, a higher pH (or higher Ca content) in the lower soil horizons can result in the breakdown of metal-humus complexes. In the lower soil layers, the organic complexing agents can be degraded by functioning microorganisms. Already established complexes in the B horizon can act as a filter, as they adsorb the traveling complexes from the upper soil horizons. A decreased water conductivity due to higher clay content can also result in the early flocculation of organo-mineral complexes.
The relocated substances can sometimes separate in the illuvial horizons. Then, organic substances are mostly enriched in the uppermost part of the illuvial horizon, whereas Fe- and Al-oxides are mostly found in the lower parts of the illuvial horizon.
Podzolization also promotes the relocation of some nutrients (Cu, Fe, Mn, Mo and P) that sometimes brings them closer to plant roots.
In different soil classification systems
The definitions in different soil classification systems are quite different. Especially soils that show pronounced other soil-forming processes in addition to podzolization are handled in different ways. The following correlations refer to soils, which have undergone advanced podzolization but lack prominent other soil-forming processes.
The term Podzols is used in the World Reference Base for Soil Resources (WRB) and in many national soil classification systems (in some of them, spelled Podsols).
The USDA soil taxonomy and the Chinese soil taxonomy call these soils Spodosols.
The Canadian system of soil classification matches Podzols with soils under the Podzolic order (e.g. Humo-Ferric Podzol).
The Australian Soil Classification uses the term Podosols.
The Brazilian Soil Classification System calls them Espodossolos.
| Physical sciences | Soil science | Earth science |
1266589 | https://en.wikipedia.org/wiki/Point%20particle | Point particle | A point particle, ideal particle or point-like particle (often spelled pointlike particle) is an idealization of particles heavily used in physics. Its defining feature is that it lacks spatial extension; being dimensionless, it does not take up space. A point particle is an appropriate representation of any object whenever its size, shape, and structure are irrelevant in a given context. For example, from far enough away, any finite-size object will look and behave as a point-like object. Point masses and point charges, discussed below, are two common cases. When a point particle has an additive property, such as mass or charge, it is often represented mathematically by a Dirac delta function. In classical mechanics there is usually no concept of rotation of point particles about their "center".
In quantum mechanics, the concept of a point particle is complicated by the Heisenberg uncertainty principle, because even an elementary particle, with no internal structure, occupies a nonzero volume. For example, the atomic orbit of an electron in the hydrogen atom occupies a volume of ~. There is nevertheless a distinction between elementary particles such as electrons or quarks, which have no known internal structure, and composite particles such as protons and neutrons, whose internal structures are made up of quarks.
Elementary particles are sometimes called "point particles" in reference to their lack of internal structure, but this is in a different sense than that discussed herein.
Point mass
Point mass (pointlike mass) is the concept, for example in classical physics, of a physical object (typically matter) that has nonzero mass, and yet explicitly and specifically is (or is being thought of or modeled as) infinitesimal (infinitely small) in its volume or linear dimensions.
In the theory of gravity, extended objects can behave as point-like even in their immediate vicinity. For example, spherical objects interacting in 3-dimensional space whose interactions are described by the Newtonian gravitation behave, as long as they do not touch each other, in such a way as if all their matter were concentrated in their centers of mass. In fact, this is true for all fields described by an inverse square law.
Point charge
Similar to point masses, in electromagnetism physicists discuss a , a point particle with a nonzero electric charge. The fundamental equation of electrostatics is Coulomb's law, which describes the electric force between two point charges. Another result, Earnshaw's theorem, states that a collection of point charges cannot be maintained in a static equilibrium configuration solely by the electrostatic interaction of the charges. The electric field associated with a classical point charge increases to infinity as the distance from the point charge decreases towards zero, which suggests that the model is no longer accurate in this limit.
In quantum mechanics
In quantum mechanics, there is a distinction between an elementary particle (also called "point particle") and a composite particle. An elementary particle, such as an electron, quark, or photon, is a particle with no known internal structure. Whereas a composite particle, such as a proton or neutron, has an internal structure.
However, neither elementary nor composite particles are spatially localized, because of the Heisenberg uncertainty principle. The particle wavepacket always occupies a nonzero volume. For example, see atomic orbital: The electron is an elementary particle, but its quantum states form three-dimensional patterns.
Nevertheless, there is good reason that an elementary particle is often called a point particle. Even if an elementary particle has a delocalized wavepacket, the wavepacket can be represented as a quantum superposition of quantum states wherein the particle is exactly localized. Moreover, the interactions of the particle can be represented as a superposition of interactions of individual states which are localized. This is not true for a composite particle, which can never be represented as a superposition of exactly-localized quantum states. It is in this sense that physicists can discuss the intrinsic "size" of a particle: The size of its internal structure, not the size of its wavepacket. The "size" of an elementary particle, in this sense, is exactly zero.
For example, for the electron, experimental evidence shows that the size of an electron is less than . This is consistent with the expected value of exactly zero. (This should not be confused with the classical electron radius, which, despite the name, is unrelated to the actual size of an electron.)
| Physical sciences | Basics_4 | Physics |
1266658 | https://en.wikipedia.org/wiki/Crystallization | Crystallization | Crystallization is the process by which solids form, where the atoms or molecules are highly organized into a structure known as a crystal. Some ways by which crystals form are precipitating from a solution, freezing, or more rarely deposition directly from a gas. Attributes of the resulting crystal depend largely on factors such as temperature, air pressure, cooling rate, and in the case of liquid crystals, time of fluid evaporation.
Crystallization occurs in two major steps. The first is nucleation, the appearance of a crystalline phase from either a supercooled liquid or a supersaturated solvent. The second step is known as crystal growth, which is the increase in the size of particles and leads to a crystal state. An important feature of this step is that loose particles form layers at the crystal's surface and lodge themselves into open inconsistencies such as pores, cracks, etc.
The majority of minerals and organic molecules crystallize easily, and the resulting crystals are generally of good quality, i.e. without visible defects. However, larger biochemical particles, like proteins, are often difficult to crystallize. The ease with which molecules will crystallize strongly depends on the intensity of either atomic forces (in the case of mineral substances), intermolecular forces (organic and biochemical substances) or intramolecular forces (biochemical substances).
Crystallization is also a chemical solid–liquid separation technique, in which mass transfer of a solute from the liquid solution to a pure solid crystalline phase occurs. In chemical engineering, crystallization occurs in a crystallizer. Crystallization is therefore related to precipitation, although the result is not amorphous or disordered, but a crystal.
Process
The crystallization process consists of two major events, nucleation and crystal growth which are driven by thermodynamic properties as well as chemical properties.
Nucleation is the step where the solute molecules or atoms dispersed in the solvent start to gather into clusters, on the microscopic scale (elevating solute concentration in a small region), that become stable under the current operating conditions. These stable clusters constitute the nuclei. Therefore, the clusters need to reach a critical size in order to become stable nuclei. Such critical size is dictated by many different factors (temperature, supersaturation, etc.). It is at the stage of nucleation that the atoms or molecules arrange in a defined and periodic manner that defines the crystal structure – note that "crystal structure" is a special term that refers to the relative arrangement of the atoms or molecules, not the macroscopic properties of the crystal (size and shape), although those are a result of the internal crystal structure.
The crystal growth is the subsequent size increase of the nuclei that succeed in achieving the critical cluster size. Crystal growth is a dynamic process occurring in equilibrium where solute molecules or atoms precipitate out of solution, and dissolve back into solution. Supersaturation is one of the driving forces of crystallization, as the solubility of a species is an equilibrium process quantified by Ksp. Depending upon the conditions, either nucleation or growth may be predominant over the other, dictating crystal size.
Many compounds have the ability to crystallize with some having different crystal structures, a phenomenon called polymorphism. Certain polymorphs may be metastable, meaning that although it is not in thermodynamic equilibrium, it is kinetically stable and requires some input of energy to initiate a transformation to the equilibrium phase. Each polymorph is in fact a different thermodynamic solid state and crystal polymorphs of the same compound exhibit different physical properties, such as dissolution rate, shape (angles between facets and facet growth rates), melting point, etc. For this reason, polymorphism is of major importance in industrial manufacture of crystalline products. Additionally, crystal phases can sometimes be interconverted by varying factors such as temperature, such as in the transformation of anatase to rutile phases of titanium dioxide.
In nature
There are many examples of natural process that involve crystallization.
Geological time scale process examples include:
Natural (mineral) crystal formation (see also gemstone);
Stalactite/stalagmite, rings formation;
Human time scale process examples include:
Snow flakes formation;
Honey crystallization (nearly all types of honey crystallize).
Methods
Crystal formation can be divided into two types, where the first type of crystals are composed of a cation and anion, also known as a salt, such as sodium acetate. The second type of crystals are composed of uncharged species, for example menthol.
Crystals can be formed by various methods, such as: cooling, evaporation, addition of a second solvent to reduce the solubility of the solute (technique known as antisolvent or drown-out), solvent layering, sublimation, changing the cation or anion, as well as other methods.
The formation of a supersaturated solution does not guarantee crystal formation, and often a seed crystal or scratching the glass is required to form nucleation sites.
A typical laboratory technique for crystal formation is to dissolve the solid in a solution in which it is partially soluble, usually at high temperatures to obtain supersaturation. The hot mixture is then filtered to remove any insoluble impurities. The filtrate is allowed to slowly cool. Crystals that form are then filtered and washed with a solvent in which they are not soluble, but is miscible with the mother liquor. The process is then repeated to increase the purity in a technique known as recrystallization.
For biological molecules in which the solvent channels continue to be present to retain the three dimensional structure intact, microbatch crystallization under oil and vapor diffusion have been the common methods.
Typical equipment
Equipment for the main industrial processes for crystallization.
Tank crystallizers. Tank crystallization is an old method still used in some specialized cases. Saturated solutions, in tank crystallization, are allowed to cool in open tanks. After a period of time the mother liquor is drained and the crystals removed. Nucleation and size of crystals are difficult to control. Typically, labor costs are very high.
Mixed-Suspension, Mixed-Product-Removal (MSMPR): MSMPR is used for much larger scale inorganic crystallization. MSMPR can crystalize solutions in a continuous manner.
Thermodynamic view
The crystallization process appears to violate the second principle of thermodynamics. Whereas most processes that yield more orderly results are achieved by applying heat, crystals usually form at lower temperaturesespecially by supercooling. However, the release of the heat of fusion during crystallization causes the entropy of the universe to increase, thus this principle remains unaltered.
The molecules within a pure, perfect crystal, when heated by an external source, will become liquid. This occurs at a sharply defined temperature (different for each type of crystal). As it liquifies, the complicated architecture of the crystal collapses. Melting occurs because the entropy (S) gain in the system by spatial randomization of the molecules has overcome the enthalpy (H) loss due to breaking the crystal packing forces:
Regarding crystals, there are no exceptions to this rule. Similarly, when the molten crystal is cooled, the molecules will return to their crystalline form once the temperature falls beyond the turning point. This is because the thermal randomization of the surroundings compensates for the loss of entropy that results from the reordering of molecules within the system. Such liquids that crystallize on cooling are the exception rather than the rule.
The nature of the crystallization process is governed by both thermodynamic and kinetic factors, which can make it highly variable and difficult to control. Factors such as impurity level, mixing regime, vessel design, and cooling profile can have a major impact on the size, number, and shape of crystals produced.
Dynamics
As mentioned above, a crystal is formed following a well-defined pattern, or structure, dictated by forces acting at the molecular level. As a consequence, during its formation process the crystal is in an environment where the solute concentration reaches a certain critical value, before changing status. Solid formation, impossible below the solubility threshold at the given temperature and pressure conditions, may then take place at a concentration higher than the theoretical solubility level. The difference between the actual value of the solute concentration at the crystallization limit and the theoretical (static) solubility threshold is called supersaturation and is a fundamental factor in crystallization.
Nucleation
Nucleation is the initiation of a phase change in a small region, such as the formation of a solid crystal from a liquid solution. It is a consequence of rapid local fluctuations on a molecular scale in a homogeneous phase that is in a state of metastable equilibrium. Total nucleation is the sum effect of two categories of nucleation – primary and secondary.
Primary nucleation
Primary nucleation is the initial formation of a crystal where there are no other crystals present or where, if there are crystals present in the system, they do not have any influence on the process. This can occur in two conditions. The first is homogeneous nucleation, which is nucleation that is not influenced in any way by solids. These solids include the walls of the crystallizer vessel and particles of any foreign substance. The second category, then, is heterogeneous nucleation. This occurs when solid particles of foreign substances cause an increase in the rate of nucleation that would otherwise not be seen without the existence of these foreign particles. Homogeneous nucleation rarely occurs in practice due to the high energy necessary to begin nucleation without a solid surface to catalyze the nucleation.
Primary nucleation (both homogeneous and heterogeneous) has been modeled as follows:
where
B is the number of nuclei formed per unit volume per unit time,
N is the number of nuclei per unit volume,
kn is a rate constant,
c is the instantaneous solute concentration,
c* is the solute concentration at saturation,
(c − c*) is also known as supersaturation,
n is an empirical exponent that can be as large as 10, but generally ranges between 3 and 4.
Secondary nucleation
Secondary nucleation is the formation of nuclei attributable to the influence of the existing microscopic crystals in the magma. More simply put, secondary nucleation is when crystal growth is initiated with contact of other existing crystals or "seeds". The first type of known secondary crystallization is attributable to fluid shear, the other due to collisions between already existing crystals with either a solid surface of the crystallizer or with other crystals themselves. Fluid-shear nucleation occurs when liquid travels across a crystal at a high speed, sweeping away nuclei that would otherwise be incorporated into a crystal, causing the swept-away nuclei to become new crystals. Contact nucleation has been found to be the most effective and common method for nucleation. The benefits include the following:
Low kinetic order and rate-proportional to supersaturation, allowing easy control without unstable operation.
Occurs at low supersaturation, where growth rate is optimal for good quality.
Low necessary energy at which crystals strike avoids the breaking of existing crystals into new crystals.
The quantitative fundamentals have already been isolated and are being incorporated into practice.
The following model, although somewhat simplified, is often used to model secondary nucleation:
where
k1 is a rate constant,
MT is the suspension density,
j is an empirical exponent that can range up to 1.5, but is generally 1,
b is an empirical exponent that can range up to 5, but is generally 2.
Growth
Once the first small crystal, the nucleus, forms it acts as a convergence point (if unstable due to supersaturation) for molecules of solute touching – or adjacent to – the crystal so that it increases its own dimension in successive layers. The pattern of growth resembles the rings of an onion, as shown in the picture, where each colour indicates the same mass of solute; this mass creates increasingly thin layers due to the increasing surface area of the growing crystal. The supersaturated solute mass the original nucleus may capture in a time unit is called the growth rate expressed in kg/(m2*h), and is a constant specific to the process. Growth rate is influenced by several physical factors, such as surface tension of solution, pressure, temperature, relative crystal velocity in the solution, Reynolds number, and so forth.
The main values to control are therefore:
Supersaturation value, as an index of the quantity of solute available for the growth of the crystal;
Total crystal surface in unit fluid mass, as an index of the capability of the solute to fix onto the crystal;
Retention time, as an index of the probability of a molecule of solute to come into contact with an existing crystal;
Flow pattern, again as an index of the probability of a molecule of solute to come into contact with an existing crystal (higher in laminar flow, lower in turbulent flow, but the reverse applies to the probability of contact).
The first value is a consequence of the physical characteristics of the solution, while the others define a difference between a well- and poorly designed crystallizer.
Size distribution
The appearance and size range of a crystalline product is extremely important in crystallization. If further processing of the crystals is desired, large crystals with uniform size are important for washing, filtering, transportation, and storage, because large crystals are easier to filter out of a solution than small crystals. Also, larger crystals have a smaller surface area to volume ratio, leading to a higher purity. This higher purity is due to less retention of mother liquor which contains impurities, and a smaller loss of yield when the crystals are washed to remove the mother liquor. In special cases, for example during drug manufacturing in the pharmaceutical industry, small crystal sizes are often desired to improve drug dissolution rate and bio-availability. The theoretical crystal size distribution can be estimated as a function of operating conditions with a fairly complicated mathematical process called population balance theory (using population balance equations).
Main crystallization processes
Some of the important factors influencing solubility are:
Concentration
Temperature
Solvent mixture composition
Polarity
Ionic strength
So one may identify two main families of crystallization processes:
Cooling crystallization
Evaporative crystallization
This division is not really clear-cut, since hybrid systems exist, where cooling is performed through evaporation, thus obtaining at the same time a concentration of the solution.
A crystallization process often referred to in chemical engineering is the fractional crystallization. This is not a different process, rather a special application of one (or both) of the above.
Cooling crystallization
Application
Most chemical compounds, dissolved in most solvents, show the so-called direct solubility that is, the solubility threshold increases with temperature.
So, whenever the conditions are favorable, crystal formation results from simply cooling the solution. Here cooling is a relative term: austenite crystals in a steel form well above 1000 °C. An example of this crystallization process is the production of Glauber's salt, a crystalline form of sodium sulfate. In the diagram, where equilibrium temperature is on the x-axis and equilibrium concentration (as mass percent of solute in saturated solution) in y-axis, it is clear that sulfate solubility quickly decreases below 32.5 °C. Assuming a saturated solution at 30 °C, by cooling it to 0 °C (note that this is possible thanks to the freezing-point depression), the precipitation of a mass of sulfate occurs corresponding to the change in solubility from 29% (equilibrium value at 30 °C) to approximately 4.5% (at 0 °C) – actually a larger crystal mass is precipitated, since sulfate entrains hydration water, and this has the side effect of increasing the final concentration.
There are limitations in the use of cooling crystallization:
Many solutes precipitate in hydrate form at low temperatures: in the previous example this is acceptable, and even useful, but it may be detrimental when, for example, the mass of water of hydration to reach a stable hydrate crystallization form is more than the available water: a single block of hydrate solute will be formed – this occurs in the case of calcium chloride);
Maximum supersaturation will take place in the coldest points. These may be the heat exchanger tubes which are sensitive to scaling, and heat exchange may be greatly reduced or discontinued;
A decrease in temperature usually implies an increase of the viscosity of a solution. Too high a viscosity may give hydraulic problems, and the laminar flow thus created may affect the crystallization dynamics.
It is not applicable to compounds having reverse solubility, a term to indicate that solubility increases with temperature decrease (an example occurs with sodium sulfate where solubility is reversed above 32.5 °C).
Cooling crystallizers
The simplest cooling crystallizers are tanks provided with a mixer for internal circulation, where temperature decrease is obtained by heat exchange with an intermediate fluid circulating in a jacket. These simple machines are used in batch processes, as in processing of pharmaceuticals and are prone to scaling. Batch processes normally provide a relatively variable quality of the product along with the batch.
The Swenson-Walker crystallizer is a model, specifically conceived by Swenson Co. around 1920, having a semicylindric horizontal hollow trough in which a hollow screw conveyor or some hollow discs, in which a refrigerating fluid is circulated, plunge during rotation on a longitudinal axis. The refrigerating fluid is sometimes also circulated in a jacket around the trough. Crystals precipitate on the cold surfaces of the screw/discs, from which they are removed by scrapers and settle on the bottom of the trough. The screw, if provided, pushes the slurry towards a discharge port.
A common practice is to cool the solutions by flash evaporation: when a liquid at a given T0 temperature is transferred in a chamber at a pressure P1 such that the liquid saturation temperature T1 at P1 is lower than T0, the liquid will release heat according to the temperature difference and a quantity of solvent, whose total latent heat of vaporization equals the difference in enthalpy. In simple words, the liquid is cooled by evaporating a part of it.
In the sugar industry, vertical cooling crystallizers are used to exhaust the molasses in the last crystallization stage downstream of vacuum pans, prior to centrifugation. The massecuite enters the crystallizers at the top, and cooling water is pumped through pipes in counterflow.
Evaporative crystallization
Another option is to obtain, at an approximately constant temperature, the precipitation of the crystals by increasing the solute concentration above the solubility threshold. To obtain this, the solute/solvent mass ratio is increased using the technique of evaporation. This process is insensitive to change in temperature (as long as hydration state remains unchanged).
All considerations on control of crystallization parameters are the same as for the cooling models.
Evaporative crystallizers
Most industrial crystallizers are of the evaporative type, such as the very large sodium chloride and sucrose units, whose production accounts for more than 50% of the total world production of crystals. The most common type is the forced circulation (FC) model (see evaporator). A pumping device (a pump or an axial flow mixer) keeps the crystal slurry in homogeneous suspension throughout the tank, including the exchange surfaces; by controlling pump flow, control of the contact time of the crystal mass with the supersaturated solution is achieved, together with reasonable velocities at the exchange surfaces. The Oslo, mentioned above, is a refining of the evaporative forced circulation crystallizer, now equipped with a large crystals settling zone to increase the retention time (usually low in the FC) and to roughly separate heavy slurry zones from clear liquid. Evaporative crystallizers tend to yield larger average crystal size and narrows the crystal size distribution curve.
DTB crystallizer
Whichever the form of the crystallizer, to achieve an effective process control it is important to control the retention time and the crystal mass, to obtain the optimum conditions in terms of crystal specific surface and the fastest possible growth. This can be achieved by a separation – to put it simply – of the crystals from the liquid mass, in order to manage the two flows in a different way. The practical way is to perform a gravity settling to be able to extract (and possibly recycle separately) the (almost) clear liquid, while managing the mass flow around the crystallizer to obtain a precise slurry density elsewhere. A typical example is the DTB (Draft Tube and Baffle) crystallizer, an idea of Richard Chisum Bennett (a Swenson engineer and later President of Swenson) at the end of the 1950s. The DTB crystallizer (see images) has an internal circulator, typically an axial flow mixer – yellow – pushing upwards in a draft tube while outside the crystallizer there is a settling area in an annulus; in it the exhaust solution moves upwards at a very low velocity, so that large crystals settle – and return to the main circulation – while only the fines, below a given grain size are extracted and eventually destroyed by increasing or decreasing temperature, thus creating additional supersaturation. A quasi-perfect control of all parameters is achieved as DTF crystallizers offer superior control over crystal size and characteristics. This crystallizer, and the derivative models (Krystal, CSC, etc.) could be the ultimate solution if not for a major limitation in the evaporative capacity, due to the limited diameter of the vapor head and the relatively low external circulation not allowing large amounts of energy to be supplied to the system.
| Physical sciences | Crystallography | null |
1266901 | https://en.wikipedia.org/wiki/Salmeterol | Salmeterol | Salmeterol is a long-acting β2 adrenergic receptor agonist (LABA) used in the maintenance and prevention of asthma symptoms and maintenance of chronic obstructive pulmonary disease (COPD) symptoms. Symptoms of bronchospasm include shortness of breath, wheezing, coughing and chest tightness. It is also used to prevent breathing difficulties during exercise (exercise-induced bronchoconstriction).
It was patented in 1983 and came into medical use in 1990. It is marketed as Serevent in the US. It is available as a dry-powder inhaler (DPI) that releases a powdered form of the drug. It was previously available as a metered-dose inhaler (MDI) but was discontinued in the US in 2002. It is available as an MDI in other countries as of 2020.
Mechanism of action
Inhaled salmeterol belongs to a group of drugs called beta-2 agonists. These drugs stimulate beta-2 receptors present in the bronchial musculature. This causes them to relax and prevent the onset and worsening of symptoms of asthma. They act on the enzyme adenyl cyclase which increases the concentration of cAMP (Cyclic adenosine monophosphate). This cyclic AMP decreases the smooth muscle tone. This drug is 10,000-times more lipid soluble than the short acting beta-2 adrenoceptor agonist, albuterol. Unlike albuterol, salmeterol becomes dissolved in the lipid bilayer of the cell membrane, and its gradual dissociation from the cell membrane provides beta-2 adrenoceptors with a supply of agonist for an extended period of time.
The primary noticeable difference of salmeterol from salbutamol, and other short-acting β2 adrenoreceptor agonists (SABAs), is its duration of action. Salmeterol lasts approximately 12 hours in comparison with salbutamol, which lasts about 4–6 hours. When used regularly every day as prescribed, inhaled salmeterol decreases the number and severity of asthma attacks. Formoterol has been demonstrated to have a faster onset of action than salmeterol as a result of a lower lipophilicity, and has also been demonstrated to be more potent—a 12 μg dose of formoterol has been demonstrated to be equivalent to a 50 μg dose of salmeterol.
Medical uses
Salmeterol is used in moderate-to-severe persistent asthma following previous treatment with a short-acting β2 adrenoreceptor agonist (SABA) such as salbutamol (albuterol).
LABAs should not be used as a monotherapy, instead, they should be used concurrently with an inhaled corticosteroid, such as beclometasone dipropionate or fluticasone propionate in the treatment of asthma to minimize serious reactions such as asthma-related deaths. Combination of inhaled corticosteroids and salmeterol (LABA) has synergistic action and reduces the frequency of asthma attacks and also makes it less severe.
In chronic obstructive pulmonary disease (COPD), LABAs may be used as monotherapy or in combination with corticosteroids. The Torch study demonstrated benefits in terms of quality of life and lung function of salmeterol alone or in combination with inhaled corticosteroids in patients with COPD
In exercise-induced bronchospasm monotherapy may be indicated in patients without persistent asthma. LABAs should not be used to treat acute symptoms.
Pregnancy and lactation
Salmeterol use during pregnancy must be decided based on the risks versus benefits to the mother. There are no well-controlled studies with salmeterol in pregnant women. Some animal studies showed developmental malformation when the mother was given several clinical doses orally. In rats, salmeterol xinafoate is excreted in the milk. However, since there is no data to show excretion of salmeterol in a mother's breast milk, a decision on whether to continue or discontinue therapy should be decided based on the important benefits it provides to the mother. Pregnant and lactating women should consult their doctors before using salmeterol.
Side effects
Due to its vasodilation properties, the common side effects of salmeterol are
dizziness,
sinus infection, and
migraine headaches.
Other side effects
muscle tremors,
hypokalemia due to direct effect on beta-2 receptors on skeletal muscle.
In most cases, salmeterol side effects are minor and either do not require treatment or can easily be treated. Certain side effects, however, should be reported to a healthcare provider immediately.
Some of these more serious side effects include
very fast heart rate,
high blood pressure, and
worsening breathing problems.
Structure-activity relationship
Salmeterol has an aryl alkyl group with a chain length of 11 atoms from the amine. This bulkiness makes the compound more lipophilic and it also makes it selective to β2 adrenergic receptors.
History
Salmeterol, first marketed and manufactured by Glaxo (now GlaxoSmithKline, GSK) in the 1980s, was released as Serevent in 1990. The product is marketed by GSK under the Allen & Hanburys brand in the UK.
In November 2005, the US Food and Drug Administration (FDA) released a health advisory, alerting the public to findings that show the use of long-acting β2 agonists could lead to a worsening of symptoms, and in some cases death.
While the use of inhaled LABAs are still recommended in asthma guidelines for the resulting improved symptom control, further concerns have been raised. A large meta-analysis of pooled results from 19 trials with 33,826 participants, suggests that salmeterol may increase the small risks of asthma-related deaths, and this additional risk is not reduced with the additional use of inhaled steroids (e.g., as with the combination product fluticasone/salmeterol).
This seems to occur because although LABAs relieve asthma symptoms, they also promote bronchial inflammation and sensitivity without warning.
Society and culture
Names
Combinations of inhaled steroids and these long-acting bronchodilators are becoming more widespread; the most common combination currently in use is fluticasone/salmeterol (brand names Seretide (UK) and Advair (US)). Another combination is budesonide/formoterol (brand name Symbicort).
| Biology and health sciences | Specific drugs | Health |
1266982 | https://en.wikipedia.org/wiki/Ranitidine | Ranitidine | {{Infobox drug
| Verifiedfields = changed
| Watchedfields = changed
| verifiedrevid = 458460406
| image = Ranitidine.svg
| image_class = skin-invert-image
| width = 250
| alt =
| image2 = File:Ranitidine-A-3D-balls.png
| alt2 =
| JAN = ranitidine hydrochloride
| pronounce =
| tradename = Zantac, others
| Drugs.com =
| MedlinePlus = a601106
| DailyMedID = Ranitidine
| pregnancy_AU = B1
| pregnancy_AU_comment =
| pregnancy_category =
| routes_of_administration = By mouth, intravenous (IV)
| class = Histamine H2 receptor antagonist, aka H2 blocker
| ATC_prefix = A02
| ATC_suffix = BA02
| ATC_supplemental = (ranitidine bismuth citrate)
| legal_AU = S4
| legal_AU_comment = / S2 (Pharmacy Medicine)
| legal_BR =
| legal_BR_comment =
| legal_CA = Rx-only
| legal_CA_comment = / OTC
| legal_DE =
| legal_DE_comment =
| legal_NZ =
| legal_NZ_comment =
| legal_UK = POM
| legal_UK_comment = / GSL/ P
| legal_US = Rx-only
| legal_US_comment = / OTC
| legal_EU = Rx-only
| legal_EU_comment =
| legal_UN =
| legal_UN_comment =
| legal_status =
| bioavailability = 50% (by mouth)
| protein_bound = 15%
| metabolism = Liver: FMOs, including FMO3; other enzymes
| metabolites =
| onset = 55–65 minutes (150 mg dose)55–115 minutes (75 mg dose)
| elimination_half-life = 2–3 hours
| duration_of_action =
| excretion = 30–70% kidney
| index2_label = HCl
| CAS_number_Ref =
| CAS_number = 66357-35-5
| CAS_number2_Ref =
| CAS_number2 = 66357-59-3
| CAS_supplemental =
| PubChem = 3001055
| PubChem2 = 3033332
| IUPHAR_ligand = 1234
| DrugBank_Ref =
| DrugBank = DB00863
| DrugBank2_Ref =
| DrugBank2 = DBSALT000487
| ChemSpiderID_Ref =
| ChemSpiderID = 4863
| ChemSpiderID2_Ref =
| ChemSpiderID2 = 43590
| UNII_Ref =
| UNII = 884KT10YB7
| UNII2_Ref =
| UNII2 = BK76465IHM
| KEGG_Ref =
| KEGG = D00422
| KEGG2_Ref =
| KEGG2 = D00673
| ChEBI_Ref =
| ChEBI = 8776
| ChEBI2_Ref =
| ChEBI2 = 8777
| ChEMBL_Ref =
| ChEMBL = 1790041
| ChEMBL2_Ref =
| ChEMBL2 = 2110372
| NIAID_ChemDB =
| PDB_ligand =
| synonyms = Dimethyl [(5-{[(2-{[1-(methylamino)-2-nitroethenyl]amino}ethyl)sulfanyl]methyl}furan-2-yl)methyl]amine
| IUPAC_name = N-(2-[(5-[(Dimethylamino)methyl]furan-2-yl)methylthio]ethyl)-''N-methyl-2-nitroethene-1,1-diamine
| C = 13 | H = 22 | N = 4 | O = 3 | S = 1
| SMILES = CNC(=C[N+](=O)[O-])NCCSCC1=CC=C(O1)CN(C)C
| StdInChI_Ref =
| StdInChI = 1S/C13H22N4O3S/c1-14-13(9-17(18)19)15-6-7-21-10-12-5-4-11(20-12)8-16(2)3/h4-5,9,14-15H,6-8,10H2,1-3H3
| StdInChI_comment =
| StdInChIKey_Ref =
| StdInChIKey = VMXUWOKSQNHOCA-UHFFFAOYSA-N
| density =
| density_notes =
| melting_point =
| melting_high =
| melting_notes =
| boiling_point =
| boiling_notes =
| solubility =
| sol_units =
| specific_rotation =
}}Ranitidine, previously sold under the brand name Zantac among others, is a medication used to decrease stomach acid production. It was commonly used in treatment of peptic ulcer disease, gastroesophageal reflux disease, and Zollinger–Ellison syndrome. It can be given by mouth, injection into a muscle, or injection into a vein.
In September 2019, the probable carcinogen N-nitrosodimethylamine (NDMA) was discovered in ranitidine products from a number of manufacturers, resulting in recalls. In April 2020, ranitidine was withdrawn from the United States market and suspended in the European Union and Australia due to these concerns.
In 2022, these concerns were confirmed in a Taiwanese nationwide population study finding "significant trends of increased liver cancer risk with an increasing dose of ranitidine" (up to 22% higher than control) and increased gastric, pancreatic, lung and overall cancer risk.
Common side effects include headaches, and pain or burning sensation if given by injection. Serious side effects may include cancer, liver problems, a slow heart rate, pneumonia, and the potential of masking stomach cancer. It is also linked to an increased risk of Clostridioides difficile colitis. Ranitidine is an H2 histamine receptor antagonist that works by blocking histamine, thus decreasing the amount of acid released by cells of the stomach.
Ranitidine was discovered in England in 1976 and came into commercial use in 1981. It is on the World Health Organization's List of Essential Medicines. It has been withdrawn at regulator request from most markets, including the United States; according to the UK NHS, it has been discontinued globally.
Former medical uses
Relief of heartburn
Short-term and maintenance therapy of gastric and duodenal ulcers
With nonsteroidal anti-inflammatory drugs (NSAIDs) to reduce the risk of ulceration Proton-pump inhibitors (PPIs) are more effective for the prevention of NSAID-induced ulcers.
Pathologic gastrointestinal (GI) hypersecretory conditions such as Zollinger–Ellison syndrome
Gastroesophageal reflux disease (GORD or GERD)
Erosive esophagitis
Part of a multidrug regimen for H. pylori eradication to minimise the risk of duodenal ulcer recurrence
Recurrent postoperative ulcer
Upper GI bleeding
For prevention of acid-aspiration pneumonitis during surgery, it can be administered preoperatively. The drug increases gastric pH, but generally has no effect on gastric volume. In a 2009 meta-analysis comparing the net benefit of PPIs and ranitidine to reduce the risk of aspiration before anaesthesia, ranitidine was found to be more effective than PPIs in reducing the volume of gastric secretions. Ranitidine may have an anti-emetic effect when administered preoperatively.
Prevention of stress-induced ulcers in critically ill patients
Used together with diphenhydramine as secondary treatment for anaphylaxis; after first-line epinephrine.
Contraindication
Ranitidine has been discontinued globally, according to the NHS, and is contraindicated due to excess cancer risk and the ready availability of H2 antagonist and PPI alternatives.
Adverse effects
These adverse effects for ranitidine have been reported as events in clinical trials:
Central nervous system
Rare reports have been made of ranitidine causing malaise, dizziness, somnolence, insomnia, and vertigo. In severely ill, elderly patients, cases of reversible mental confusion, agitation, depression, and hallucinations have been reported.
Cardiovascular
Arrhythmias such as tachycardia, bradycardia, atrioventricular block, and premature ventricular beats have also been reported.
Gastrointestinal
All drugs in the H2 receptor blocker class of medicines have the potential to cause vitamin B12 deficiency, secondary to a reduction in food-bound vitamin B12 absorption. Elderly patients taking H2 receptor antagonists are more likely to require B12 supplementation than those not taking such drugs. H2 blockers may also reduce the absorption of drugs (azole antifungals, calcium carbonate) that require an acidic stomach. In addition, multiple studies suggest the use of H2 receptor antagonists such as ranitidine may increase the risk of infectious diarrhoea, including traveller's diarrhoea and salmonellosis. A 2005 study found that by suppressing acid-mediated breakdown of proteins, ranitidine may lead to an elevated risk of developing food or drug allergies, due to undigested proteins then passing into the GI tract, where sensitisation occurs. Patients who take these agents develop higher levels of immunoglobulin E against food, whether they had prior antibodies or not. Even months after discontinuation, an elevated level of IgE in 6% of patients was still found in the study.
Liver
Cholestatic hepatitis, liver failure, hepatitis, and jaundice have been noted, and require immediate discontinuation of the drug. Blood tests can reveal an increase in liver enzymes or eosinophilia, although in rare instances, severe cases of hepatotoxicity may require a liver biopsy.
Lungs
Ranitidine and other histamine H2 receptor antagonists may increase the risk of pneumonia in hospitalised patients. Ranitidine increases the risk of community-acquired pneumonia in adults and children.
Blood
Thrombocytopenia is a rare but known side effect. Drug-induced thrombocytopenia usually takes weeks or months to appear, but may appear within 12 hours of drug intake in a sensitised individual. Typically, the platelet count falls to 80% of normal, and thrombocytopenia may be associated with neutropenia and anemia.
Skin
Rash, including rare cases of erythema multiforme, and rare cases of hair loss and vasculitis have been seen.
Precautions
Disease-related concerns
Relief of symptoms due to the use of ranitidine does not exclude the presence of a gastric malignancy. In addition, with kidney or liver impairment, ranitidine must be used with caution. It should be avoided in patients with porphyria, as it may precipitate an attack.
Children
In children, the use of gastric acid inhibitors has been associated with an increased risk for development of acute gastroenteritis and community-acquired pneumonia. A cohort analysis including over 11,000 neonates reported an association of H2 blocker use, and an increased incidence of necrotizing enterocolitis in very-low-birth-weight (VLBW) neonates. In addition, about a six-fold increase in mortality, necrotizing enterocolitis, and infection such as sepsis, pneumonia, urinary tract infection was reported in patients receiving ranitidine in a cohort analysis of 274 VLBW neonates.
Drug tests
Ranitidine may return a false positive result with some commercial urine drug screening kits for testing for drugs of abuse.
Cancer-causing due to inherent instability
In June 2019, Valisure informed the US Food and Drug Administration (FDA) that Zantac-branded and generic ranitidine resulted in very high levels of NDMA in the human body "due to an inherent instability of the ranitidine molecule".
In September 2019, the FDA acknowledged that ranitidine medicines, including some products sold under the brand name Zantac, contained a nitrosamine impurity called N-nitrosodimethylamine (NDMA), classified as a probable human carcinogen, at unacceptable levels. Health Canada announced that it was assessing NDMA in ranitidine and requested that manufacturers stop the distribution of ranitidine products in Canada until the NDMA levels in the products are found to be safe. Health Canada announced that ranitidine drugs were being recalled by Sandoz Canada, Apotex Inc., Pro Doc Limitée, Sanis Health Inc., and Sivem Pharmaceuticals ULC. The European Medicines Agency (EMA) started a European Union-wide review of ranitidine medicines at the request of the European Commission.
In October 2019, the US. FDA observed that the third-party laboratory that found very high levels of NDMA was using higher temperatures in its tests to detect nitrosamine impurities. The NDMA was mostly generated by the added heat, but the higher temperatures are recommended for using a gas chromatography–mass spectrometry method to test for NDMA in valsartan and angiotensin II receptor blockers. The FDA stated that it recommends using a Liquid Chromatography-High Resolution Mass Spectrometry (LC-HRMS) testing protocol to test samples of ranitidine. Its LC-HRMS testing method does not use elevated temperatures, and has shown the presence of much lower levels of NDMA in ranitidine medicines than were reported by the third-party laboratory. International regulators using similar LC-MS testing methods have also shown the presence of lower but still unacceptable levels of NDMA in ranitidine samples. The FDA provided additional guidance about using another LC-MS method based on a triple-quadrupole MS platform.
In September 2019, Sandoz issued a "precautionary distribution stop" of all medicines containing ranitidine, followed a few days later by a recall of ranitidine hydrochloride capsules in the United States. The Italian Medicines Agency recalled all ranitidine that uses an active pharmaceutical ingredient from Saraca Laboratories. The Federal Union of German Associations of Pharmacists (Arzneimittelkommission der Deutschen Apotheker) published a list of recalled products, as did the Therapeutic Goods Administration in Australia.
In November 2019, the FDA stated that its tests found levels of NDMA in ranitidine and nizatidine that are similar to those that one may typically ingest with common foods such as grilled or smoked meats. The FDA also stated that its simulated gastric fluid model tests and simulated intestinal fluid model tests indicated that NDMA is not formed when exposed to acid in the stomach with a normal diet. The FDA advised companies to recall their ranitidine if testing shows levels of NDMA above the acceptable daily intake (96 nanograms per day or 0.32 parts per million for ranitidine). At the same time, it indicated that some levels of NDMA found in medicines still exceeded the agency's acceptable levels.
In December 2019, the FDA asked manufacturers of ranitidine and nizatidine products to expand their NDMA testing to include all lots of the medication before making them available to consumers.
By the end of 2019, ranitidine had already fallen from the 40th most commonly prescribed medication in the United States in 2018, to 53rd place for 2019, with about 13.6million prescriptions for the year, versus nearly 19 million the previous year. This reflects total prescriptions for all of 2019, while safety concerns affected sales in only the final 4 months of the year.
In April 2020, new FDA testing and evaluation prompted by information from third-party laboratories confirmed that NDMA levels increase in ranitidine even under normal storage conditions, and NDMA has been found to increase significantly in samples stored at higher temperatures, including those at which the product may be exposed during distribution and handling by consumers. The testing also showed that the level of NDMA increases as ranitidine medication ages. These conditions may raise the NDMA level above the acceptable daily intake limit. The EMA completed and issued their EU-wide review at the end of the month and the European Commission suspended all ranitidine products in the EU.
In August 2020, the EMA provided guidance to marketing authorization holders for avoiding the presence of nitrosamine impurities and asked them to review all chemical and biological human medicines for the presence of nitrosamines and to test the products at risk. In September 2020, the FDA issued guidance about the control of nitrosamine impurities in human drugs. An implementation plan was issued in February 2021.
In 2022, these concerns were confirmed in a Taiwanese nationwide population study finding "significant trends of increased liver cancer risk with an increasing dose of ranitidine" (22% higher than control) and increased gastric, pancreatic, lung (26%, 35%, and 17% respectively), but "only liver cancer displayed a significant association with long-term ranitidine use" and "there was no continuous dose–response relationship among the other individual cancers". Overall cancer risk also increased by 10% (p < 0.001).
The FDA issued revised guidelines about nitrosamine impurities in September 2024.
List of recalls
In September 2019, Apotex recalled all over-the-counter ranitidine tablets sold in the United States at Walmart, Rite Aid, and Walgreens. These retailers, along with CVS, removed Zantac and some generics from their shelves.
In October 2019, the Medicines and Healthcare products Regulatory Agency of the United Kingdom (UK) issued a drug alert for ranitidine "... to proactively communicate the recall to hospitals, pharmacies, dispensing practices, retailers and wholesalers in the UK." This included all Zantac-branded preparations, along with all generic preparations of ranitidine from Teva UK Limited, Rosemont Pharmaceuticals Limited, Omega Pharma Limited and Galpharm International Limited, Perrigo Company plc, Creo Pharma Limited and Tillomed Laboratories Limited, OTC Concepts Ltd, Relonchem Ltd, Noumed Life Sciences Ltd, and Medreich Plc., Accord Healthcare, Medley Pharma Limited, and Medreich Plc.
In October 2019, the Department of Health and Social Care of the United Kingdom issued a supply distribution alert (SDA/2019/005) for all oral formulations of ranitidine.
In October 2019, Sanofi recalled all over-the-counter Zantac in the United States and Canada, Perrigo issued a worldwide recall of ranitidine, Dr. Reddy's issued a recall of all ranitidine products in the United States, and Novitium Pharma recalled all ranitidine hydrochloride capsules in the US.
In November 2019, Aurobindo Pharma, Amneal Pharmaceuticals, American Health Packaging, Golden State Medical Supply, and Precision Dose recalled some lots of ranitidine tablets, capsules, and syrup.
In December 2019, Glenmark Pharmaceutical Inc., USA, recalled some lots of ranitidine tablets.
In January 2020, Appco Pharma LLC and Northwind Pharmaceuticals recalled some lots of ranitidine tablets and capsules.
In February 2020, American Health Packaging recalled some lots of ranitidine tablets manufactured by Amneal Pharmaceuticals.
In April 2020, the FDA requested a manufacturer's market withdrawal of ranitidine, meaning that ranitidine products would not be available for prescription or over-the-counter sale in the US.
In April 2020, the Committee for Medicinal Products for Human Use (CHMP) of the European Medicines Agency recommended the suspension of all ranitidine medicines in the European Union because of the presence of unacceptable levels of NDMA. Text was copied from this source which is copyright European Medicines Agency. Reproduction is authorized provided the source is acknowledged.
A ranitidine manufacturer requested a re-examination of the decision, but in December 2020, the EMA confirmed its recommendation to suspend all ranitidine medicines in the European Union. The UK National Health Service (NHS) Web site said "Ranitidine is not currently available in the UK or globally... It's not yet known whether it will be available again in future." A March 2024 review left the message up.
In 2021, Solara Active Pharma Sciences, which supplies ranitidine active pharmaceutical ingredient (API), said that it had mitigated the risks of the formation of NDMA during the manufacturing of ranitidine API. The company was granted a revised certificate by the European Directorate for the Quality of Medicines and Healthcare, which proves that the API complies with certain European rules. GlaxoSmithKline, Sanofi, and Teva said they had no plans to reintroduce the drug in the EU, but Accord Healthcare considered the possible reintroduction of ranitidine. However a control strategy regarding NDMA formation through the end of the product’s shelf life, despite heat, time and digestion due to endogenous formation from the API, would be required.
Pharmacology
Mechanism of action
Ranitidine is a competitive, reversible inhibitor of the action of histamine at the histamine H2 receptors found in gastric parietal cells. This results in decreased gastric acid secretion and gastric volume, and reduced hydrogen ion concentration. Ranitidine's acid-lowering effect is more pronounced for basal and nocturnal acid secretion than it is for food-stimulated acid secretion. Additional indirect effects of ranitidine are decreased pepsin secretion and increased nitrate-reducing bacterial flora.
Pharmacokinetics
Oral absorption: 50%
Protein binding: 15%
Metabolism: The major metabolite in the urine is ranitidine N-oxide, which represents less than 4% of the dose. Other metabolites of ranitidine include ranitidine S-oxide (1%) and desmethyl ranitidine (1%).
Half-life elimination: With normal renal function, ranitidine taken orally has a half-life of 2.5–3.0 hours. If taken intravenously, the half-life is generally 2.0–2.5 hours in a patient with normal kidney function and normal creatinine clearance. In patients with kidney dysfunction, the half-life may increase to 4 to 5 hours.
Excretion: The primary route of excretion is the urine. In addition, about 30% of the orally administered dose is collected in the urine as unabsorbed drug in 24 hours.
Elderly
In the elderly population, the plasma half-life of ranitidine is prolonged to 3–4 hours secondary to decreased kidney function causing decreased clearance.
Children
In general, studies of pediatric patients (aged one month to 16 years) have shown no significant differences in pharmacokinetic parameter values in comparison to healthy adults, when correction is made for body weight.
History
Ranitidine was first prepared in England as AH19065 by John Bradshaw in the summer of 1977 in the Ware research laboratories of Allen and Hanburys, part of the former Glaxo organisation. Its development was a response to the first in class histamine H2 receptor antagonist, cimetidine, developed by Sir James Black at Smith, Kline and French, and launched in the United Kingdom as Tagamet''' in November 1976. Both companies eventually merged as GlaxoSmithKline (GSK), following a sequence of mergers and acquisitions, starting with the integration of Allen and Hanbury's Ltd and Glaxo to form Glaxo Group Research in 1979, and ultimately with the merger of Glaxo Wellcome and SmithKline Beecham in 2000. Ranitidine was the result of a rational drug-design process using what was by then a fairly refined model of the histamine H2 receptor and quantitative structure-activity relationships.
Glaxo refined the model further, by replacing the imidazole ring of cimetidine with a furan ring with a nitrogen-containing substituent, and in doing so developed ranitidine. Ranitidine was found to have a far-improved tolerability profile (i.e. fewer adverse drug reactions), longer-lasting action, and 10 times the activity of cimetidine. Ranitidine has 10% of the affinity that cimetidine has to CYP450, so it causes fewer side effects, but other H2 blockers famotidine and nizatidine have no CYP450 significant interactions.
Ranitidine was introduced in 1981, and was the world's biggest-selling prescription drug by 1987. Subsequently, it was largely superseded by the more effective proton-pump inhibitor (PPI) class of drugs, with omeprazole becoming the biggest-selling drug for many years. When omeprazole and ranitidine were compared in a study of 144 people with severe inflammation and erosions or ulcers of the oesophagus, 85% of those treated with omeprazole healed within eight weeks, compared with 50% of those given ranitidine. In addition, the omeprazole group reported earlier relief of heartburn symptoms.
In September 2019, the probable carcinogen N-nitrosodimethylamine (NDMA) was discovered in ranitidine products from a number of manufacturers, resulting in recalls; in April 2020, it was withdrawn from the United States market and suspended in Europe and Australia.
Preparations
Preparations of ranitidine products include oral tablets (75, 150, and 300 mg), effervescent tablets, and syrups, and injectable solutions; with doses of specific ranitidine product preparations are available over-the-counter (OTC) in various countries. In the United Kingdom, only the lowest-strength, 75-mg tablet was available to purchase without a prescription. In Australia, packs containing seven or 14 doses of the 150-mg tablet were available in supermarkets, small packs of 150-mg and 300-mg tablets were schedule 2 pharmacy medicines. Larger doses and pack sizes required a prescription. In the United States, 75- and 150-mg tablets were available OTC. In India, it is sold under multiple brand names.
| Biology and health sciences | Antihistamines | Health |
22387642 | https://en.wikipedia.org/wiki/Common%20envelope | Common envelope | In astronomy, a common envelope (CE) is gas that contains a binary star system. The gas does not rotate at the same rate as the embedded binary system. A system with such a configuration is said to be in a common envelope phase or undergoing common envelope evolution.
During a common envelope phase the embedded binary system is subject to drag forces from the envelope which cause the separation of the two stars to decrease. The phase ends either when the envelope is ejected to leave the binary system with much smaller orbital separation, or when the two stars become sufficiently close to merge and form a single star. A common envelope phase is short-lived relative to the lifetime of the stars involved.
Evolution through a common envelope phase with ejection of the envelope can lead to the formation of a binary system composed of a compact object with a close companion. Cataclysmic variables, X-ray binaries and systems of close double white dwarfs or neutron stars are examples of systems of this type which can be explained as having undergone common envelope evolution. In all these examples there is a compact remnant (a white dwarf, neutron star or black hole) which must have been the core of a star which was much larger than the current orbital separation. If these systems have undergone common envelope evolution then their present close separation is explained. Short-period systems containing compact objects are sources of gravitational waves and Type Ia supernovae.
Predictions of the outcome of common envelope evolution are uncertain.
A common envelope is sometimes confused with a contact binary. In a common envelope binary system the envelope does not generally rotate at the same rate as the embedded binary system; thus it is not constrained by the equipotential surface passing through the L2 Lagrangian point. In a contact binary system the shared envelope rotates with the binary system and fills an equipotential surface.
Formation
A common envelope is formed in a binary star system when the orbital separation decreases rapidly or one of the stars expands rapidly.
The donor star will start mass transfer when it overfills its Roche lobe and as a consequence the orbit will shrink further causing it to overflow the Roche lobe even more, which accelerates the mass transfer, causing the orbit to shrink even faster and the donor to expand more. This leads to the run-away process of dynamically unstable mass transfer. In some case the receiving star is unable to accept all material, which leads to the formation of a common envelope engulfing the companion star.
Evolution
The donor's core does not participate in the expansion of the stellar envelope and the formation of the common envelope, and the common envelope will contain two objects: the core of the original donor and the companion star. These two objects (initially) continue their orbital motion inside the common envelope. However, it is thought that because of drag forces inside the gaseous envelope, the two objects lose energy, which brings them in a closer orbit and actually increases their orbital velocities. The loss of orbital energy is assumed to heat up and expand the envelope, and the whole common-envelope phase ends when either the envelope is expelled into space, or the two objects inside the envelope merge and no more energy is available to expand or even expel the envelope. This phase of the shrinking of the orbit inside the common envelope is known as a spiral-in.
Observational manifestations
Common envelope events (CEEs) are difficult to observe. Their existence has been mainly inferred indirectly from presence in the Galaxy of binary systems that can not be explained by any other mechanism. Observationally CEEs should be brighter than typical novae but fainter than typical supernovae. The photosphere of the common envelope should be relatively cool—at about 5,000 K—emitting a red spectrum. However its large size should lead to a large luminosity—on the order of that of a red supergiant. A common envelope event should begin with a sharp rise in luminosity followed by a few months long plateau of constant luminosity (much like that of type II-P supernova) powered by the recombination of hydrogen in the envelope. After that the luminosity should decrease rapidly.
Several events that resemble the description above have been observed in past. These events are called luminous red novae (LRNe). They are subset of a broader class of events called intermediate-luminosity red transients (ILRTs). They have relatively slow expansion velocities of 200–1000 km/s and total radiated energies are 1038 to 1040 J.
The possible CEEs that have been observed so far include:
M85 OT2006-1, possible ejection of the whole envelope.
V1309 Scorpii, a possible star merger.
M31 RV
V838 Monocerotis
Ou 5, a planetary nebula whose progenitor was a common envelope binary
| Physical sciences | Stellar astronomy | Astronomy |
20888255 | https://en.wikipedia.org/wiki/Isopropyl%20alcohol | Isopropyl alcohol | Isopropyl alcohol (IUPAC name propan-2-ol and also called isopropanol or 2-propanol) is a colorless, flammable, organic compound with a pungent alcoholic odor.
Isopropyl alcohol, an organic polar molecule, is miscible in water, ethanol, and chloroform, demonstrating its ability to dissolve a wide range of substances including ethyl cellulose, polyvinyl butyral, oils, alkaloids, and natural resins. Notably, it is not miscible with salt solutions and can be separated by adding sodium chloride in a process known as salting out. It forms an azeotrope with water, resulting in a boiling point of 80.37 °C and is characterized by its slightly bitter taste. Isopropyl alcohol becomes viscous at lower temperatures, freezing at −89.5 °C, and has significant ultraviolet-visible absorbance at 205 nm. Chemically, it can be oxidized to acetone or undergo various reactions to form compounds like isopropoxides or aluminium isopropoxide. As an isopropyl group linked to a hydroxyl group (chemical formula ) it is the simplest example of a secondary alcohol, where the alcohol carbon atom is attached to two other carbon atoms. It is a structural isomer of propan-1-ol and ethyl methyl ether. They all have the formula .
It was first synthesized in 1853 by Alexander William Williamson and later produced for cordite preparation. It is produced through hydration of propene or hydrogenation of acetone, with modern processes achieving anhydrous alcohol through azeotropic distillation. Beyond its production, isopropyl alcohol serves in medical settings as a rubbing alcohol and hand sanitizer, and in industrial and household applications as a solvent. It is a common ingredient in products such as antiseptics, disinfectants and detergents. More than a million tonnes are produced worldwide annually. Despite its utility, isopropyl alcohol poses safety risks due to its flammability and potential for peroxide formation. Its ingestion or absorption leads to toxic effects including central nervous system depression and coma, primarily treated through supportive measures.
Properties
Isopropyl alcohol is miscible in water, ethanol, and chloroform, as it is an organic polar molecule. It dissolves ethyl cellulose, polyvinyl butyral, many oils, alkaloids, and natural resins. Unlike ethanol or methanol, isopropyl alcohol is not miscible with salt solutions and can be separated from aqueous solutions by adding a salt such as sodium chloride. The process is colloquially called salting out, and causes concentrated isopropyl alcohol to separate into a distinct layer.
Isopropyl alcohol forms an azeotrope with water, which gives a boiling point of and a composition of 87.7% by mass (91% by volume) isopropyl alcohol. It has a slightly bitter taste, and is not safe to drink.
Isopropyl alcohol becomes increasingly viscous with decreasing temperature and freezes at . Mixtures with water have higher freezing points: 99% at , 91% (the azeotrope) at , and 70% at .
Isopropyl alcohol has a maximal absorbance at 205 nm in an ultraviolet-visible spectrum.
Reactions
Isopropyl alcohol can be oxidized to acetone, which is the corresponding ketone. This can be achieved using oxidizing agents such as chromic acid, or by dehydrogenation of isopropyl alcohol over a heated copper catalyst:
Isopropyl alcohol is often used as both solvent and hydride source in the Meerwein-Ponndorf-Verley reduction and other transfer hydrogenation reactions. Isopropyl alcohol may be converted to 2-bromopropane using phosphorus tribromide, or dehydrated to propene by heating with sulfuric acid.
Like most alcohols, isopropyl alcohol reacts with active metals such as potassium to form alkoxides that are called isopropoxides. With titanium tetrachloride, isopropyl alcohol reacts to give titanium isopropoxide:
This and similar reactions are often conducted in the presence of base.
The reaction with aluminium is initiated by a trace of mercury to give aluminium isopropoxide.
History
Isopropyl alcohol was first synthesized by the chemist Alexander William Williamson in 1853. He achieved this by heating a mixture of propene and sulfuric acid.
Standard Oil produced isopropyl alcohol by hydrating propene. Isopropyl alcohol was oxidized to acetone for the preparation of cordite, a smokeless, low explosive propellant.
Production
In 1994, 1.5 million tonnes of isopropyl alcohol were produced in the United States, Europe, and Japan. It is primarily produced by combining water and propene in a hydration reaction or by hydrogenating acetone. There are two routes for the hydration process and both processes require that the isopropyl alcohol be separated from water and other by-products by distillation. Isopropyl alcohol and water form an azeotrope, and simple distillation gives a material that is 87.9% by mass isopropyl alcohol and 12.1% by mass water. Pure (anhydrous) isopropyl alcohol is made by azeotropic distillation of the wet isopropyl alcohol using either diisopropyl ether or cyclohexane as azeotroping agents.
Biological
Small amounts of isopropyl alcohol are produced in the body in diabetic ketoacidosis.
Indirect hydration
Indirect hydration reacts propene with sulfuric acid to form a mixture of sulfate esters. This process can use low-quality propene, and is predominant in the USA. These processes give primarily isopropyl alcohol rather than 1-propanol, because adding water or sulfuric acid to propene follows Markovnikov's rule. Subsequent hydrolysis of these esters by steam produces isopropyl alcohol, by distillation. Diisopropyl ether is a significant by-product of this process; it is recycled back to the process and hydrolyzed to give the desired product.
Direct hydration
Direct hydration reacts propene and water, either in gas or liquid phase, at high pressures in the presence of solid or supported acidic catalysts. This type of process usually requires higher-purity propylene (> 90%). Direct hydration is more commonly used in Europe.
Hydrogenation of acetone
Isopropyl alcohol can be prepared via the hydrogenation of acetone, but this approach involves an extra step compared to the above methods, as acetone is itself normally prepared from propene via the cumene process. cost is primarily driven by raw material cost, and this way is economical when acetone is cheaper than propylene as a byproduct of phenol production (the coexistence of two ways on most markets allows them to balance the prices).
A known issue is the formation of MIBK and other self-condensation products. Raney nickel was one of the original industrial catalysts, modern catalysts are often supported bimetallic materials.
Uses
In 1990, 45,000 metric tonnes of isopropyl alcohol were used in the United States, mostly as a solvent for coatings or for industrial processes. In that year, 5400 metric tonnes were used for household purposes and in personal care products. Isopropyl alcohol is popular in particular for pharmaceutical applications, due to its low toxicity. Some isopropyl alcohol is used as a chemical intermediate. Isopropyl alcohol may be converted to acetone, but the cumene process is more significant.
Solvent
Isopropyl alcohol dissolves a wide range of non-polar compounds. It evaporates quickly and the typically available grades tend to not leave behind oil traces when used as a cleaning fluid unlike some other common solvents. It is also relatively non-toxic. Thus, it is used widely as a solvent and as a cleaning fluid, especially where there are oils or oil based residues which are not easily cleaned with water, conveniently evaporating and (depending on water content and other variables) posing less of a risk of corrosion or rusting than plain water. Together with ethanol, n-butanol, and methanol, it belongs to the group of alcohol solvents.
Isopropyl alcohol is commonly used for cleaning eyeglasses, electrical contacts, audio or video tape heads, DVD and other optical disc lenses, bongs, and for removing thermal paste from heatsinks on CPUs and other IC packages.
Intermediate
Isopropyl alcohol is esterified to give isopropyl acetate, another solvent. It reacts with carbon disulfide and sodium hydroxide to give sodium isopropylxanthate, which has use as an herbicide and an ore flotation reagent. Isopropyl alcohol reacts with titanium tetrachloride and aluminium metal to give titanium and aluminium isopropoxides, respectively, the former a catalyst, and the latter a chemical reagent. This compound may serve as a chemical reagent in itself, by acting as a dihydrogen donor in transfer hydrogenation.
Medical
Rubbing alcohol, hand sanitizer, and disinfecting pads typically contain a 60–70% solution of isopropyl alcohol or ethanol in water. Water is required to open up membrane pores of bacteria, which acts as a gateway for isopropyl alcohol. A 75% solution in water may be used as a hand sanitizer. Isopropyl alcohol is used as a water-drying aid for the prevention of otitis externa, better known as swimmer's ear.
Inhaled isopropyl alcohol can be used for treating nausea in some settings by placing a disinfecting pad under the nose.
Early uses as an anesthetic
Although isopropyl alcohol can be used for anesthesia, its many negative attributes or drawbacks prohibit this use. Isopropyl alcohol can also be used similarly to ether as a solvent or as an anesthetic by inhaling the fumes or orally. Early uses included using the solvent as general anesthetic for small mammals and rodents by scientists and some veterinarians. However, it was soon discontinued, as many complications arose, including respiratory irritation, internal bleeding, and visual and hearing problems. In rare cases, respiratory failure leading to death in animals was observed.
Automotive
Isopropyl alcohol is a major ingredient in "gas dryer" fuel additives. In significant quantities, water is a problem in fuel tanks, as it separates from gasoline and can freeze in the supply lines at low temperatures. Alcohol does not remove water from gasoline, but the alcohol solubilizes water in gasoline. Once soluble, water does not pose the same risk as insoluble water, as it no longer accumulates in the supply lines and freezes but is dissolved within the fuel itself. Isopropyl alcohol is often sold in aerosol cans as a windshield or door lock deicer. Isopropyl alcohol is also used to remove brake fluid traces from hydraulic braking systems, so that the brake fluid (usually DOT 3, DOT 4, or mineral oil) does not contaminate the brake pads and cause poor braking. Mixtures of isopropyl alcohol and water are also commonly used in homemade windshield washer fluid.
Laboratory
As a biological specimen preservative, isopropyl alcohol provides a comparatively non-toxic alternative to formaldehyde and other synthetic preservatives. Isopropyl alcohol solutions of 70–99% are used to preserve specimens.
Isopropyl alcohol is often used in DNA extraction. A lab worker adds it to a DNA solution to precipitate the DNA, which then forms a pellet after centrifugation. This is possible because DNA is insoluble in isopropyl alcohol.
Semiconductors
Isopropyl alcohol is used as an additive in alkaline anisotropic etching of monocrystalline silicon, such as with potassium hydroxide or tetramethylammonium hydroxide. This process is used in texturing of silicon solar cells and microfabrication (e.g. in MEMS devices). Isopropyl alcohol increases the anisotropy of the etch by increasing the etch rate of [100] plane relative to higher indexed planes.
Safety
Isopropyl alcohol vapor is denser than air and is flammable, with a flammability range of between 2% and 12.7% in air. It should be kept away from heat, sparks, and open flame. Distillation of isopropyl alcohol over magnesium has been reported to form peroxides, which may explode upon concentration. Isopropyl alcohol can react with air and oxygen over time to form unstable peroxides that can explode.
Toxicology
Isopropyl alcohol, via its metabolites, is somewhat more toxic than ethanol, but considerably less toxic than ethylene glycol or methanol. Death from ingestion or absorption of even relatively large quantities is rare. Both isopropyl alcohol and its metabolite, acetone, act as central nervous system (CNS) depressants. Poisoning can occur from ingestion, inhalation, or skin absorption. Symptoms of isopropyl alcohol poisoning include flushing, headache, dizziness, CNS depression, nausea, vomiting, anesthesia, hypothermia, low blood pressure, shock, respiratory depression, and coma. Overdoses may cause a fruity odor on the breath as a result of its metabolism to acetone.
Isopropyl alcohol does not cause an anion gap acidosis, but it produces an osmolal gap between the calculated and measured osmolalities of serum, as do the other alcohols. The findings of acetone without acidosis leads to the sine qua non of "ketosis without acidosis."
Isopropyl alcohol is oxidized to form acetone by alcohol dehydrogenase in the liver and has a biological half-life in humans between 2.5 and 8.0 hours. Unlike methanol or ethylene glycol poisoning, the metabolites of isopropyl alcohol are considerably less toxic, and treatment is largely supportive. Furthermore, there is no indication for the use of fomepizole, an alcohol dehydrogenase inhibitor, unless co-ingestion with methanol or ethylene glycol is suspected.
In forensic pathology, people who have died as a result of diabetic ketoacidosis or alcoholic ketoacidosis, with no isopropyl alcohol ingestion, usually have detectable blood concentrations of isopropyl alcohol of 1 to 40 mg/dL, while those by fatal isopropyl alcohol ingestion usually have blood concentrations of hundreds of mg/dL.
| Physical sciences | Alcohols | Chemistry |
2619735 | https://en.wikipedia.org/wiki/Flow%20%28mathematics%29 | Flow (mathematics) | In mathematics, a flow formalizes the idea of the motion of particles in a fluid. Flows are ubiquitous in science, including engineering and physics. The notion of flow is basic to the study of ordinary differential equations. Informally, a flow may be viewed as a continuous motion of points over time. More formally, a flow is a group action of the real numbers on a set.
The idea of a vector flow, that is, the flow determined by a vector field, occurs in the areas of differential topology, Riemannian geometry and Lie groups. Specific examples of vector flows include the geodesic flow, the Hamiltonian flow, the Ricci flow, the mean curvature flow, and Anosov flows. Flows may also be defined for systems of random variables and stochastic processes, and occur in the study of ergodic dynamical systems. The most celebrated of these is perhaps the Bernoulli flow.
Formal definition
A flow on a set is a group action of the additive group of real numbers on . More explicitly, a flow is a mapping
such that, for all and all real numbers and ,
It is customary to write instead of , so that the equations above can be expressed as (the identity function) and (group law). Then, for all the mapping is a bijection with inverse This follows from the above definition, and the real parameter may be taken as a generalized functional power, as in function iteration.
Flows are usually required to be compatible with structures furnished on the set . In particular, if is equipped with a topology, then is usually required to be continuous. If is equipped with a differentiable structure, then is usually required to be differentiable. In these cases the flow forms a one-parameter group of homeomorphisms and diffeomorphisms, respectively.
In certain situations one might also consider s, which are defined only in some subset
called the of . This is often the case with the flows of vector fields.
Alternative notations
It is very common in many fields, including engineering, physics and the study of differential equations, to use a notation that makes the flow implicit. Thus, is written for and one might say that the variable depends on the time and the initial condition . Examples are given below.
In the case of a flow of a vector field on a smooth manifold , the flow is often denoted in such a way that its generator is made explicit. For example,
Orbits
Given in , the set is called the orbit of under . Informally, it may be regarded as the trajectory of a particle that was initially positioned at . If the flow is generated by a vector field, then its orbits are the images of its integral curves.
Examples
Algebraic equation
Let be a time-dependent trajectory which is a bijective function. Then a flow can be defined by
Autonomous systems of ordinary differential equations
Let be a (time-independent) vector field
and the solution of the initial value problem
Then is the flow of the vector field . It is a well-defined local flow provided that the vector field
is Lipschitz-continuous. Then is also Lipschitz-continuous wherever defined. In general it may be hard to show that the flow is globally defined, but one simple criterion is that the vector field is compactly supported.
Time-dependent ordinary differential equations
In the case of time-dependent vector fields , one denotes where is the solution of
Then is the time-dependent flow of . It is not a "flow" by the definition above, but it can easily be seen as one by rearranging its arguments. Namely, the mapping
indeed satisfies the group law for the last variable:
One can see time-dependent flows of vector fields as special cases of time-independent ones by the following trick. Define
Then is the solution of the "time-independent" initial value problem
if and only if is the solution of the original time-dependent initial value problem. Furthermore, then the mapping is exactly the flow of the "time-independent" vector field .
Flows of vector fields on manifolds
The flows of time-independent and time-dependent vector fields are defined on smooth manifolds exactly as they are defined on the Euclidean space and their local behavior is the same. However, the global topological structure of a smooth manifold is strongly manifest in what kind of global vector fields it can support, and flows of vector fields on smooth manifolds are indeed an important tool in differential topology. The bulk of studies in dynamical systems are conducted on smooth manifolds, which are thought of as "parameter spaces" in applications.
Formally: Let be a differentiable manifold. Let denote the tangent space of a point Let be the complete tangent manifold; that is, Let
be a time-dependent vector field on ; that is, is a smooth map such that for each and , one has that is, the map maps each point to an element of its own tangent space. For a suitable interval containing 0, the flow of is a function that satisfies
Solutions of heat equation
Let be a subdomain (bounded or not) of (with an integer). Denote by its boundary (assumed smooth).
Consider the following heat equation on , for ,
with the following initial value condition in .
The equation on corresponds to the Homogeneous Dirichlet boundary condition. The mathematical setting for this problem can be the semigroup approach. To use this tool, we introduce the unbounded operator defined on by its domain
(see the classical Sobolev spaces with
and
is the closure of the infinitely differentiable functions with compact support in for the norm).
For any , we have
With this operator, the heat equation becomes and . Thus, the flow corresponding to this equation is (see notations above)
where is the (analytic) semigroup generated by .
Solutions of wave equation
Again, let be a subdomain (bounded or not) of (with an integer). We denote by its boundary (assumed smooth).
Consider the following wave equation on (for ),
with the following initial condition in and
Using the same semigroup approach as in the case of the Heat Equation above. We write the wave equation as a first order in time partial differential equation by introducing the following unbounded operator,
with domain on (the operator is defined in the previous example).
We introduce the column vectors
(where and ) and
With these notions, the Wave Equation becomes and .
Thus, the flow corresponding to this equation is
where is the (unitary) semigroup generated by
Bernoulli flow
Ergodic dynamical systems, that is, systems exhibiting randomness, exhibit flows as well. The most celebrated of these is perhaps the Bernoulli flow. The Ornstein isomorphism theorem states that, for any given entropy , there exists a flow , called the Bernoulli flow, such that the flow at time , i.e. , is a Bernoulli shift.
Furthermore, this flow is unique, up to a constant rescaling of time. That is, if , is another flow with the same entropy, then , for some constant . The notion of uniqueness and isomorphism here is that of the isomorphism of dynamical systems. Many dynamical systems, including Sinai's billiards and Anosov flows are isomorphic to Bernoulli shifts.
| Mathematics | Dynamical systems | null |
25219833 | https://en.wikipedia.org/wiki/Shared%20transport | Shared transport | Shared transport or shared mobility is a transportation system where travelers share a vehicle either simultaneously as a group (e.g. ride-sharing) or over time (e.g. carsharing or bike sharing) as personal rental, and in the process share the cost of the journey. It is a transportation strategy that allows users to access transportation services on an as-needed basis, and can be regarded as a hybrid between private vehicle use and mass or public transport. Shared mobility is an umbrella term that encompasses a variety of transportation modes including carsharing, Bicycle-sharing systems, ridesharing companies, carpools, and microtransit.
Each shared mobility service has unique attributes that have a range of impacts on travel behavior, the environment, and the development of cities and urban areas. Some impacts of shared mobility include enhanced transportation accessibility as well as reduced driving and decreased personal vehicle ownership.
Shared mobility programs often yield a variety of environmental, social, and transportation system benefits. These are primarily related to personal vehicle usage and ownership, and vehicle miles or kilometers traveled (VMT/VKT). Shared mobility networks also retain the potential to expand the reach of public transportation by addressed gaps in existing public transportation systems. They can also provide economic benefits to users in the form of cost savings in some cases.
Shared transport systems include carsharing (also called car clubs in the UK), bicycle sharing (also known as PBS or public bicycle systems), carpools and vanpools (aka ride-sharing or lift-sharing), real-time ridesharing, slugging (casual carpooling), community buses and vans, demand responsive transit (DRT), paratransit, a range of taxi projects and even hitchhiking and its numerous variants.
Shared transport is taking on increasing importance as a key strategy for reducing greenhouse gas and other emissions from the transport sector in the face of the global climate emergency by finding ways of getting more intensive use of vehicles on the road. Together with other emerging automotive technologies such as vehicle electrification, connected vehicles and autonomous driving, shared transports form a future mobility vision called Connected, Autonomous, Shared and Electric (CASE) Mobility.
A somewhat different form of shared transport is the "shared taxi", a vehicle which follows a predetermined route and takes anybody waiting for it, more like a bus than a taxi.
History
Shared mobility is a subgroup of the larger sharing economy. The sharing economy is a term that encompasses a wide variety of services, usually involving the online transactions of goods or services as part of a peer-to-peer marketplace. Innovations in social networking, location-based services, and Internet technologies have enabled shared mobility to develop and expand rapidly. By improving efficiency, providing cost savings, and monetizing underused resources, shared mobility services have become widely used in many cities around the world. Although the proliferation of tech-enabled shared mobility has occurred mostly within the last decade, shared mobility services are not a new phenomenon. The first carsharing program was established in 1948 in Zurich, Switzerland, and the first bikesharing program began in 1965 in Amsterdam, Netherlands.
Smartphone applications and location data have increased the feasibility of shared transportation services, including carsharing companies and mobile app-based vehicle for hire companies.
Auto rickshaws
Auto rickshaws carry people and goods in many developing countries. Also known as a three-wheeler, Samosa, tempo, tuk-tuk, trishaw, auto, rickshaw, autorick, bajaj, rick, tricycle, mototaxi, baby taxi or lapa in popular parlance, they are motorized version of the traditional pulled rickshaw or cycle rickshaw. They are an essential form of urban transport, both as vehicles for hire and for private use, in many developing countries, and a form of novelty transport in many Eastern countries.
Bikesharing
Bicycle-sharing systems allow users to access and use a shared fleet of bicycles, typically located within a given spatial boundary. These systems are mostly concentrated in cities or other urban areas and bikes or stations are normally unattended and always accessible. This availability during most or all of the day makes bikesharing an on-demand mobility option.
The first bikesharing system debuted in Amsterdam in 1965, under the name ‘White Bikes.’ The bicycles were left unlocked around the city to be used by anyone in need of transportation. Bikesharing systems have since exploded in popularity starting in the mid-2000s due to advances in information technology (IT) that have improved bikesharing communications and tracking. As of April 2016, there were 99 U.S. cities with technology-enabled public bikesharing systems, with approximately 32,200 bikes and 3,400 stations.
Three major types of bikesharing systems have emerged: public bikesharing (docked and dockless/free floating), closed campus bikesharing and peer-to-peer (P2P) bikesharing. Most bikesharing systems are public and allow anyone to access a bicycle for a fee, typically in daily, monthly or annual membership fees. Public bikesharing programs can be station-based (docked), or dockless (also known as free floating). Dockless systems are deployed within a geo-fenced area. Dockless systems were first introduced in Germany in the early 2000s via the Call a Bike program. Major bikesharing operators in North America include: Motivate, Social Bicycles, Spin, ofo, Mobike, and LimeBike. E-bikesharing systems (or Pedlec) have also been growing in popularity, particularly in Europe. Social Bicycles began testing an e-bikesharing program, called Jump, in San Francisco in Summer 2017.
Studies have been conducted that analyze bikesharing impacts on modal shift. A 2014 UC Berkeley study suggests that in larger cities, bikesharing programs remove riders from crowded or high-use bus transit systems. In smaller cities, bikesharing improves access from bus lines, filling in gaps in the public transit system. In addition, those living in larger cities report decreased rail usage as a result of increased cost savings and reduced travel times. The study also found that half of the bikesharing members surveyed reduced their personal vehicle usage due to bikesharing.
Carsharing
Carsharing refers to a model of vehicle sharing where users access cars on an as-needed basis, and often pay by time of reservation or miles driven. As of January 2015, there were 23 carsharing operators in the U.S. amounting to over 1.1 million members and over 16,000 vehicles. As of January 2017, there were 39 carsharing organizations in North America serving 1.9 million members with a collective fleet of 24,629 vehicles. (these numbers do not include P2P carsharing; they include roundtrip and one-way carsharing operations.).
Roundtrip carsharing
Roundtrip carsharing is one of the earliest carsharing models, granting members access to a shared vehicle fleet. As the name suggests, roundtrip carsharing requires users to return to the same location where they accessed the vehicle.
One of the largest North American-based roundtrip carsharing operators is Zipcar, which operates more than 12,000 vehicles in urban areas on college campuses and at airports across the United States, Canada, the United Kingdom, Spain, France, Belgium, Turkey and Taiwan. There have been numerous studies that document behavioral changes associated with roundtrip carsharing programs. A 2004 study on City CarShare in San Francisco, CA found that nearly 30% of members reduced car ownership by one of more cars; two-thirds of members reported that they opted not to purchase an additional vehicle. This reduced car ownership typically translates into reduced driving, and thus lowered energy consumption and greenhouse gas emissions. Carsharing programs also affect usage patterns of other travel modes. A 2011 study by UC Berkeley researchers found that roundtrip carsharing has a mixed impact on public transit and non-motorized modal use, with the same proportion of respondents increasing and decreasing usage of these modes. The impact on carpooling and non-motorized transportation, however, was found to be positive. The same study documented a 27% to 43% reduction in vehicle miles traveled and a 34% to 41% reduction in greenhouse gas emissions among households due to roundtrip carsharing.
One-way carsharing
One-way carsharing varies from roundtrip carsharing in that it grants members more flexibility in pickup and dropoff location. In one-way carsharing—also known as point-to-point carsharing—members can access a vehicle at one location and end their trip in another location. As of September 2015, companies that offered one-way functionality in the U.S. include car2go, GIG, ReachNow, Zipcar, and BlueIndy. As of January 2015, about 35% of North American carsharing fleets were one-way capable.
A 2016 study of one-way carsharing operator, car2go, in five North American cities found that 2% to 5% of members sold a vehicle, and 7% to 10% postponed a vehicle purchase due to their carsharing membership. Moreover, estimated VMT impacts due to carsharing ranged from −6% to −16% per car2go household, and GHG emissions changed by −4% to −18%.
Personal vehicle sharing (PVS) and P2P carsharing
Personal vehicle sharing (PVS) is a carsharing service model that allows short-term access to privately owned vehicles. P2P carsharing, a subset of PVS, employs privately owned vehicles made available for shared usage by members of a P2P member base. P2P carsharing companies differ from other carsharing operators in that users provide the free-floating vehicle fleet using their personally owned vehicles. P2P carsharing operators in North America include Getaround and Turo (formerly RelayRides), and as of May 2015, there were eight active P2P operators in North America.
One 2014 study found that the top three reasons for using P2P carsharing are convenience and availability, monetary savings, and expanded mobility options. Another study documented that personal vehicle sharing services can expand the geographic range of vehicle sharing services by renting underused autos and therefore lowering vehicle usage requirements. However, fear of sharing personal assets was cited as one of the primary barriers to the adoption of P2P sharing services.
Ridesharing
Ridesharing services enable shared rides between drivers and passengers that have similar origins and destination pairings. Ridesharing includes both vanpooling and carpooling. Vanpooling involves a grouping of between seven and 15 people traveling in a van, and carpooling refers to groups of less than seven people traveling together in one vehicle.
Ridesharing is distinct from ridesourcing (or TNCs), like Uber and Lyft in that the driver typically decides trip origin, destination, and any deviations to accommodate one or more additional passengers. Drivers and riders have the same origin, destination, or potentially share multiple proximate destinations, with a common purpose of conserving resources, saving money, or saving time. Driver earnings from ridesharing are regulated in the U.S. by the Internal Revenue Service, and as of January 2017, they were capped at 53.5 cents per mile for business travel by car.
Both technology-enabled ridesharing organizations and more informal ridesharing programs exist. Examples of technology-enabled ridesharing companies are BlaBlaCar, Carma Carpool, Scoop, and Waze Carpool. These services usually require participants to join either through a membership, website, or mobile application. Potential drivers and riders can then post driving routes or preferred travel routes and the ridesharing service will connect riders with passengers that share similar origin destination pairings. More informal ad hoc ridesharing programs include slugging (also known as casual carpooling). Casual carpooling is an informal form of commuter ridesharing operating in Washington, D.C.; Houston, Texas; and San Francisco, California.
Casual carpooling has been in existence for over 30 years, is entirely run informally by its users, and does not use a mobile application or information communication technology. In one study in the San Francisco Bay Area in 2014, researchers interviewed, observed, and surveyed participants at multiple casual carpooling locations. The study found that motivations for casual carpooling participation include: convenience, time savings, and monetary savings, while environmental and community-based motivations ranked low. Casual carpooling is an efficient transportation option for these commuters, while environmental sustainability benefits are a positive byproduct. Seventy-five percent of casual carpool users were previously public transit riders, and over 10% formerly drove alone.
In the U.S., the modal share of ridesharing has declined since the 1970s. In 1970, The U.S. Census found that about 20% of American workers commuted to work by carpool. The American Community Survey has found that the carpooling modal share has declined to around nine percent as of 2013, though it still remains the second most popular mode of travel in the U.S., next to driving alone.
On-demand ride services
On-demand ride services include ridehailing, ridesplitting, and E-hail for taxis. They are services that provide rides on-demand, usually in passenger cars, for a fee.
On-demand vehicle for hire
"Transportation network company" is a regulatory classification coined by the California Public Utilities Commission in 2013, and it has been subsequently used by other U.S. states to refer to services like Lyft and Uber. These include point-to-point on-demand rides, typically hailed, coordinated, and paid for via smartphone and from drivers using their own personal vehicles. Transportation experts have called these services "ridesourcing" or "ridehailing" to distinguish these services from ridesharing and to clarify that drivers do not share a destination with their passengers. Ridehailing companies have spread around the world and include: Uber, Lyft, Ola Cabs, DiDi, Grab, Gett, Cabify, Careem, Easy Taxi, and Fasten, among others. As of August 2017, 2 million people drive for Uber every week.
These companies have faced criticism for adversely impacting traffic congestion, the environment, and public safety. A study of ridehailing users in San Francisco in 2014 evaluated modal shifts due to ridehailing and found that, if ridehailing were unavailable, 39% of respondents would have taken a taxi and 33% would have used a form of public transit. Four percent entered a public transit station as their origin or destination, suggesting ridehailing may serve as a first-/last-mile trip to or from public transit in some cases. Another study of ridehailing users in Denver and Boulder, Colorado found that a third of respondents would have taken public transit, biked, or walked instead of using a ridehailing service. Another third would have driven in a personal vehicle, and 12% would not have made the trip. These city-specific differences suggest that travel behavior impacts due to these services could be dependent on location. Only New York City and San Francisco have studied the vehicle miles traveled (VMT) implications of ridehailing services. Both studies found that Uber and Lyft are increasing VMT, with the heaviest impacts seen in some of the busiest areas of each city. However, both of these studies do not take into consideration modal shift changes.
Ridesplitting
Ridesplitting involves splitting both a ride and fare in a vehicle with others traveling in the same general direction. These services allow dynamic matching and route variation in real time as passengers request pickups. The user cost of ridesplitting services is lower than the cost of regular ridesourcing services, since the riders are sharing one ride and splitting the associated costs. Yet, ridesplitting may lead to detour and inconvenience effects for the users. Ridesplitting services are generally only available as an option in cities with denser and more established ridesourcing markets. Ridesplitting is even less studied than ridesourcing, and therefore travel behavior impacts are not yet well understood.
E-Hail services
E-Hail services are a mode of transportation by which taxis can be reserved via Internet or mobile phone applications maintained by either a third-party provider or the taxi company. Examples of e-Hail services include Curb, Flywheel, Arro, Hailo, and iTaxi. In response to competition from ridesourcing companies, e-Hail taxi services have experienced rapid growth. As of October 2014, 80% of San Francisco taxis reported using Flywheel, an e-Hail app. As of February 2015, Flywheel was active in six cities, and Curb was active in about 60 U.S. cities. Since they use taxis, e-Hail services charge local taxi rates and do not use demand-based pricing during periods of higher ride demand, as ridesourcing services often do.
Microtransit
Microtransit is a technology-enabled private transit service that often uses shuttles or vans and is characterized by flexible scheduling, flexible routing, or both. Current microtransit operators include Chariot (acquired by Ford in September 2016) and Via. Defunct operators include Bridj and Leap Transit. Two forms of microtransit have emerged, including fixed-route with fixed-schedule services and flexible-route with on-demand scheduling. Chariot, which started in San Francisco and now operates in Austin, New York, and Seattle, functions similarly to public transit and runs 15-seater vans along pre-determined routes. Chariot determines new routes by “crowdsourcing” potential customer demand and then launching a new route once enough demand is indicated. Via is an example of flexible route, on-demand microtransit and currently operates in New York City, Washington DC, and Chicago. In New York City, users request a ride using Via's app and a shared van will pick them up with other travelers heading in a similar direction. The service is dynamic, without static routes, and shifts routes based on expected traffic and rider demand. Via charges a fare of $5 to $7 per ride in New York City, depending on the method of booking. Both Chariot and Via conform to the IRS “transit pass” standard, allowing them to qualify for pre-tax commuter benefits.
Microtransit services have also gained interest among some public transit operators, who see the technology as an opportunity to provide higher quality or more flexible public transit services to their users. In some instances, like the (now defunct) RideKC: Bridj pilot project in Kansas City, Missouri, public-private partnerships have been formed to provide microtransit services. The RideKC: Bridj pilot ultimately ceased operations as it failed to attract enough riders, with only nine percent of riders taking more than 10 trips. The lack of a targeted marketing campaign, relatively high vehicle ownership rates in Kansas City, and low existing public transit mode share in the city were possible reasons for the low ridership of the pilot project. Public-private microtransit partnerships have the potential to improve service and increase public transit ridership, but steps must be taken to appropriately evaluate demand for the service before launching.
Courier network services
Courier Network Services (CNS) provide delivery services of packages, food, and other items for compensation using their own transportation and are connected with shippers and customers through an online app or platform. In P2P delivery services, someone who signs up and is approved by the platform can use their own vehicle or bicycle to conduct a delivery. There are many business models within P2P delivery services. Postmates couriers make deliveries using their own bicycles, scooters, or cars. They charge a delivery fee plus a service charge of nine percent of the value of the goods being delivered. Instacart delivers groceries for a $4 to $10 fee, depending on how long the delivery takes to complete. The proliferation of these services, where couriers use their personally owned vehicles or bicycles, could reduce the need for delivery companies to maintain and own their own delivery fleet.
Some CNS models that have emerged also incorporate on-demand ride services (e.g., TNCs) that deliver packages. CNS deliveries are either made in separate trips or in multiple-purpose trips that may also serve passengers simultaneously. Sidecar and Uber have incorporated passenger, food delivery, and package delivery services.
Scootersharing
Scootersharing is a recent application of the sharing economy within the transportation space. Scootersharing companies took inspiration from the fourth generation bikesharing strategy, but replaced bicycles with GPS-tracked electric scooters. These scooters are also “dockless”, and are dropped off and picked up from any location within an urban area.
Due to the lower speed of scooters and their electric assistance, it is easier for commuters to use them and for companies to invest in a fleet of them. Many scootersharing companies have been founded in the past few years. This includes Bird, Lime, Bolt, Skip, Scoot Networks, and Spin.
Because of its growing popularity, some cities have also looked to ban certain scootersharing companies, taking on similar strategies to ridesharing bans. In San Francisco, the city created a Powered Scooter Share Permit Program that limits the number of companies that could operate scooters, and the amount of scooters. Cities that enforced similar regulations cite how scooters are more commonly ridden on sidewalks instead of bike lanes and could injure pedestrians. Other reasons would also be the lack of these companies enforcing riders to wear safety gear such as helmets.
Compared to the other forms of shared mobility, scootersharing can be more hyper-localized and can hypothetically better address the last mile problem. Because scootersharing does not have much market adoption right now because it is a new form of transportation, there are no academic studies that can effectively measure its impact. Overall, it provides urban mobility with fewer carbon emissions compared to automobiles. They take up less space than bikes, so they have potential to increase transit ridership to and from bus lines.
Enabling technology: smartphone apps
Smartphones represent one of the most important transportation innovations of the 21st century. A variety of factors are changing the way people think about mobility including: demographic shifts, advancements in geo-spatial routing and computing power, the use of cloud technologies, faster wireless networks capable of carrying greater bandwidth, congestion, and heightened awareness about the environment and climate change. Mobility consumers are increasingly using smartphone applications, dubbed “apps” for an array of transportation use cases. More people are starting their trips with smartphones to plan routes, seek departure information for the next bus or railcar, find a taxi via an e-Hail app, or source a private driver through services, such as Lyft or Uber. Some factors driving transportation app growth are time savings; financial savings; incentives; and gamification.
Future of shared mobility and automation
Self-driving cars, in conjunction with shared mobility, have the potential to greatly increase the viability and user base of shared transportation services in the future. There has been great interest in the idea of shared automated vehicle (SAV) services in recent years. This interest is likely due to the highly publicized AV development space, as well as the popularity of ridesourcing services and the realization that operating cost per mile of mobility services may substantially decrease compared to current prices, with automation. Many experts, companies, public agencies, and universities are at the initial stages of exploring the potential impacts of SAVs.
A few pilots have launched involving ridesourcing services and automated vehicles. Uber began testing an AV service open to frequent customers in Pittsburgh, Pennsylvania in September 2016. Waymo (formerly Google's Self-Driving Car Project) has been testing an autonomous vehicle ride service in Arizona. Also during September 2016 in Singapore, nuTonomy and Grab partnered to offer a similar AV ridesourcing service in a business district called “One North.” These SAV services require an engineer to closely monitor the system at all times. There have also been several automated shuttle service pilots around the world, although all are in the initial testing phase and operate in a low-speed setting. Low-speed SAV shuttle companies include: EasyMile EZ10, Local Motors, Auro, and Navya SAS.
The impact that SAV services may have on travel behavior, other transportation modes, the environment, and cities in general remains uncertain. As real-world deployment of SAVs has been extremely limited, most studies on the subject develop or modify existing models of travel behavior and include SAVs, with assumptions regarding their operations and vehicle types. Studies predict a modal shift away from private vehicle trips due to SAVs under certain sharing scenarios. The impact SAV services may have on VMT and congestion is uncertain as well, with some studies predicting that roadway capacity may be freed up due to more efficient operations and right-sizing of vehicles.
Shared transport modes
Real-time ridesharing
Flexible carpooling
Green travel
Hail and ride
High-occupancy vehicle lanes
Hitchhiking
Slugging
The commons
Truck sharing
Vanpooling
Group travel
Carpooling (ride sharing)
Demand responsive transport
Paratransit
Share taxi
Rental travel
Bicycle-sharing system
Carsharing
| Technology | Motorized road transport | null |
6292760 | https://en.wikipedia.org/wiki/Bed%20%28geology%29 | Bed (geology) | In geology, a bed is a layer of sediment, sedimentary rock, or volcanic rock "bounded above and below by more or less well-defined bedding surfaces".
A bedding surface is three-dimensional surface, planar or curved, that visibly separates each successive bed (of the same or different lithology) from the preceding or following bed. Where bedding surfaces occur as cross-sections, e.g., in a 2-dimensional vertical cliff face of horizontal strata, are often referred to as bedding contacts. Within conformable successions, each bedding surface acted as the depositional surface for the accumulation of younger sediment.
Definitions
Specifically in sedimentology, a bed can be defined in one of two major ways. First, Campbell and Reineck and Singh use the term bed to refer to a thickness-independent layer comprising a coherent layer of sedimentary rock, sediment, or pyroclastic material bounded above and below by surfaces known as bedding planes. By this definition of bed, laminae are small beds that constitute the smallest (visible) layers of a hierarchical succession and often, but not always, internally comprise a bed.
Alternatively, a bed can be defined by thickness where a bed is a coherent layer of sedimentary rock, sediment, or pyroclastic material greater than 1 cm thick and a lamina is a coherent layer of sedimentary rock, sediment, or pyroclastic material less than 1 cm thick. This method of defining bed versus lamina is frequently used in textbooks, e.g., Collinson & Mountney or Miall. Both definitions have merit and the choice of which one to use will depend on the focus of the specific study on a case by case basis.
Interpretation
Typically, but not always, bedding surfaces record changes in either the rate or type of accumulating sediment that created the underlying bed. Typically, they represent either a period of nondeposition, erosional truncation, shift in flow or sediment regime, abrupt change in composition, or combination of these as a result of changes in environmental conditions. As a result, a bed is typically, but not always, interpreted to represent a single period of time when sediments or pyroclastic material accumulated during uniform and steady paleoenvironmental conditions. However, some bedding surfaces may be postdepositional features either formed or enhanced by diagenetic processes or weathering.
The relationship between bedding surfaces controls the gross geometry of a bed. Most commonly, the bottom and top surfaces of beds are subparallel to parallel to each other. However, some bedding surfaces of a bed are nonparallel, e.g., wavy, or curved. Differing combinations of nonparallel bedding surfaces results in beds of widely varying geometric shapes such as uniform-tabular, tabular-lenticular, curved-tabular, wedge-shaped, and irregular beds.
Types
Types of beds include cross-beds and graded beds. Cross-beds, or "sets," are not layered horizontally and are formed by a combination of local deposition on the inclined surfaces of ripples or dunes, and local erosion. Graded beds show a gradual change in grain or clast sizes from one side of the bed to the other. A normal grading occurs where there are larger grain sizes on the older side, while an inverse grading occurs where there are smaller grain sizes on the older side.
Bed thickness
Bed thickness is a basic and important characteristic of beds. Besides mapping stratigraphic units and interpreting sedimentary facies, the analysis of bed thickness can be used to recognize breaks in sedimentation, cyclic sedimentation patterns, and gradual environmental changes. Such sedimentological studies are typically based on the hypothesis that the thicknesses of stratigraphic units follows a log-normal distribution. Differing nomenclatures for the bed and laminae thickness have been proposed by various authors, including McKee and Weir, Ingram, and Reineck and Singh. However, none of them have been universally accepted by Earth scientists. In the practice of engineering geology, a standardized nomenclature is used for describing bed thickness in Australia, the European Union, and the United Kingdom.
Examples of widely used bed thickness classifications include Tucker (1982) and McKee and Weir (1953).
Bed in lithostratigraphy
According to both the North American Stratigraphic Code and International Stratigraphic Guide, a bed is the smallest formal lithostratigraphic unit that can be used for sedimentary rocks. A bed, a stratum, is the smallest formal unit in the hierarchy of sedimentary lithostratigraphic units and is lithologically distinguishable from other layers above and below. Customarily, only distinctive beds, i.e. key beds, marker beds, that are particularly useful for stratigraphic purposes are given proper names and considered formal lithostratigraphic units.
In case of volcanic rocks, the lithostratigraphic unit equivalent to a bed is a flow. A flow is “...a discrete, extrusive, volcanic rock body distinguishable by texture, composition, order of superposition, paleomagnetism, or other objective criteria.” A flow is a part of a member as a bed of sedimentary rock is a part of a member.
Engineering considerations
In geotechnical engineering a bedding surface often forms a discontinuity that may have a large influence on the mechanical behaviour (strength, deformation, etc.) of soil and rock masses in tunnel, foundation, or slope construction.
Geologic principles
These are the principles which apply to all geologic features, and can be used to describe the order of events in a feature's geologic history.
The Law of Superposition states that younger rocks are deposited above older rocks, and remain that way as long as the beds have not been overturned through tectonic activities. This is used to date the stratigraphy and their relative ages.
The Law of Original Horizontality states that beds are deposited horizontally due to gravity. If the beds are not horizontal, then that is an indication that they have been tilted or warped by geologic processes.
The Law of Lateral Continuity states that the bed deposits extend laterally in all directions. This implies that two places separated by erosional features with similar rocks may have originally been continuous.
The law of Cross-Cutting Relationships states that any feature which cuts through another is the younger of the two. This can include faults or igneous dikes cutting through sedimentary bedding.
| Physical sciences | Sedimentology | Earth science |
6294249 | https://en.wikipedia.org/wiki/Type-I%20superconductor | Type-I superconductor | The interior of a bulk superconductor cannot be penetrated by a weak magnetic field, a phenomenon known as the Meissner effect. When the applied magnetic field becomes too large, superconductivity breaks down. Superconductors can be divided into two types according to how this breakdown occurs. In type-I superconductors, superconductivity is abruptly destroyed via a first order phase transition when the strength of the applied field rises above a critical value Hc. This type of superconductivity is normally exhibited by pure metals, e.g. aluminium, lead, and mercury. The only alloys known up to now which exhibit type I superconductivity are tantalum silicide (TaSi2). and BeAu
The covalent superconductor SiC:B, silicon carbide heavily doped with boron, is also type-I.
Depending on the demagnetization factor, one may obtain an intermediate state. This state, first described by Lev Landau, is a phase separation into macroscopic non-superconducting and superconducting domains forming a Husimi Q representation.
This behavior is different from type-II superconductors which exhibit two critical magnetic fields. The first, lower critical field occurs when magnetic flux vortices penetrate the material but the material remains superconducting outside of these microscopic vortices. When the vortex density becomes too large, the entire material becomes non-superconducting; this corresponds to the second, higher critical field.
The ratio of the London penetration depth λ to the superconducting coherence length ξ determines whether a superconductor is type-I or type-II. Type-I superconductors are those with , and type-II superconductors are those with .
| Physical sciences | Electrical circuits | Physics |
802913 | https://en.wikipedia.org/wiki/Semiconductor%20detector | Semiconductor detector | A semiconductor detector in ionizing radiation detection physics is a device that uses a semiconductor (usually silicon or germanium) to measure the effect of incident charged particles or photons.
Semiconductor detectors find broad application for radiation protection, gamma and X-ray spectrometry, and as particle detectors.
Detection mechanism
In semiconductor detectors, ionizing radiation is measured by the number of charge carriers set free in the detector material which is arranged between two electrodes, by the radiation. Ionizing radiation produces free electrons and electron holes. The number of electron-hole pairs is proportional to the energy of the radiation to the semiconductor. As a result, a number of electrons are transferred from the valence band to the conduction band, and an equal number of holes are created in the valence band. Under the influence of an electric field, electrons and holes travel to the electrodes, where they result in a pulse that can be measured in an outer circuit, as described by the Shockley-Ramo theorem. The holes travel in the opposite direction and can also be measured. As the amount of energy required to create an electron-hole pair is known, and is independent of the energy of the incident radiation, measuring the number of electron-hole pairs allows the energy of the incident radiation to be determined.
The energy required to produce electron-hole-pairs is very low compared to the energy required to produce paired ions in a gas detector. Consequently, in semiconductor detectors the statistical variation of the pulse height is smaller and the energy resolution is higher. As the electrons travel fast, the time resolution is also very good, and is dependent upon rise time. Compared with gaseous ionization detectors, the density of a semiconductor detector is very high, and charged particles of high energy can give off their energy in a semiconductor of relatively small dimensions.
Detector types
Silicon detectors
Most silicon particle detectors work, in principle, by doping narrow (usually around 100 micrometers wide) silicon strips to turn them into diodes, which are then reverse biased. As charged particles pass through these strips, they cause small ionization currents that can be detected and measured. Arranging thousands of these detectors around a collision point in a particle accelerator can yield an accurate picture of what paths particles take. Silicon detectors have a much higher resolution in tracking charged particles than older technologies such as cloud chambers or wire chambers. The drawback is that silicon detectors are much more expensive than these older technologies and require sophisticated cooling to reduce leakage currents (noise source). They also suffer degradation over time from radiation, however, this can be greatly reduced thanks to the Lazarus effect.
Diamond detectors
Diamond detectors have many similarities with silicon detectors but are expected to offer significant advantages – in particular a high radiation hardness and very low drift currents. They are also suited to neutron detection. At present, however, they are much more expensive and more difficult to manufacture.
Germanium detectors
Germanium detectors are mostly used for gamma spectroscopy in nuclear physics, as well as x-ray spectroscopy. While silicon detectors cannot be thicker than a few millimeters, germanium can have a sensitive layer (depletion region) thickness of centimeters, and therefore can be used as a total absorption detector for gamma rays up to a few MeV. These detectors are also called high-purity germanium detectors (HPGe) or hyperpure germanium detectors. Before current purification techniques were refined, germanium crystals could not be produced with purity sufficient to enable their use as spectroscopy detectors. Impurities in the crystals trap electrons and holes, ruining the performance of the detectors. Consequently, germanium crystals were doped with lithium ions (Ge(Li)), in order to produce an intrinsic region in which the electrons and holes would be able to reach the contacts and produce a signal.
When germanium detectors were first developed, only very small crystals were available. Low efficiency was the result, and germanium detector efficiency is still often quoted in relative terms to a "standard" 3″ x 3″ NaI(Tl) scintillation detector. Crystal growth techniques have since improved, allowing detectors to be manufactured that are as large as or larger than commonly available NaI crystals, although such detectors cost more than €100,000 (US$113,000).
, HPGe detectors commonly use lithium diffusion to make an n+ ohmic contact, and boron implantation to make a p+ contact. Coaxial detectors with a central n+ contact are referred to as n-type detectors, while p-type detectors have a p+ central contact. The thickness of these contacts represents a dead layer around the surface of the crystal within which energy depositions do not result in detector signals. The central contact in these detectors is opposite to the surface contact, making the dead layer in n-type detectors smaller than the dead layer in p-type detectors. Typical dead layer thicknesses are several hundred micrometers for a Li diffusion layer and a few tenths of a micrometer for a B implantation layer.
The major drawback of germanium detectors is that they must be cooled to liquid nitrogen temperatures to produce spectroscopic data. At higher temperatures, the electrons can easily cross the band gap in the crystal and reach the conduction band, where they are free to respond to the electric field, producing too much electrical noise to be useful as a spectrometer. Cooling to liquid nitrogen temperature (77K) reduces thermal excitations of valence electrons so that only a gamma ray interaction can give an electron the energy necessary to cross the band gap and reach the conduction band. Cooling with liquid nitrogen is inconvenient, as the detector requires hours to cool down to operating temperature before it can be used, and cannot be allowed to warm up during use. Ge(Li) crystals could never be allowed to warm up, as the lithium would drift out of the crystal, ruining the detector. HPGe detectors can be allowed to warm up to room temperature when not in use.
Commercial systems became available that use advanced refrigeration techniques (for example pulse tube refrigerator) to eliminate the need for liquid nitrogen cooling.
Germanium detectors with multi-strip electrodes, orthogonal on opposing faces, can indicate the 2-D location of the ionization trail within a large single crystal of Ge. Detectors like this have been used in COSI balloon-born astronomy missions (NASA, 2016) and will be used in an orbital observatory (NASA, 2025) Compton Spectrometer and Imager (COSI).
Because germanium detectors are highly efficient in photon detection, they can be used for a variety of additional applications. High-purity germanium detectors are used by Homeland Security to differentiate between naturally occurring radioactive material (NORM) and weaponized or otherwise harmful radioactive material. They are also used in monitering the environment due to the concern of the use of nuclear power. Finally, high-purity germanium detectors are used for medical imaging and nuclear physics research, making them a rather diverse detector as far as applications go.
Cadmium telluride and cadmium zinc telluride detectors
Cadmium telluride (CdTe) and cadmium zinc telluride (CZT) detectors have been developed for use in X-ray spectroscopy and gamma spectroscopy. The high density of these materials means they can effectively attenuate X-rays and gamma-rays with energies of greater than 20 keV that traditional silicon-based sensors are unable to detect. The wide band gap of these materials also means they have high resistivity and are able to operate at, or close to, room temperature (~295K) unlike germanium-based sensors. These detector materials can be used to produce sensors with different electrode structures for imaging and high-resolution spectroscopy. However, CZT detectors are generally unable to match the resolution of germanium detectors, with some of this difference being attributable to poor positive charge-carrier transport to the electrode. Efforts to mitigate this effect have included the development of novel electrodes to negate the need for both polarities of carriers to be collected.
Integrated Systems
Semiconductor detectors are often commercially integrated into larger systems for various radiation measurement applications.
Automated Sample Changing for Germanium Detectors
Gamma spectrometers using HPGe detectors are often used for measurement of low levels of gamma-emitting radionuclides in environmental samples, which requires a low background environment, usually achieved by enclosing the sample and detector in a lead shield known as a 'lead castle'. Automated systems have been developed to sequentially move a number of samples into and out of the lead castle for measurement. Due to the complexities of opening the shield and moving the samples, this automation has traditionally been expensive, but lower-cost autosamplers have recently been introduced.
Radioactive Waste Assay Machines
Semiconductor detectors especially HPGe are often integrated into devices for characterising packaged radioactive waste. This can be as simple as detectors being mounted on a moveable platform to be brought to an area for in-situ measurements and paired with shielding to restrict the field-of-view of the detector to the area of interest for one-shot "open detector geometry" measurements, or for waste in drums, systems such as the Segmented Gamma Scanner (SGS) combine a semiconductor detector with integrated mechatronics to rotate the item and scan the detector across different sections. If the detector field of view is scanned across small areas of the item in multiple axes as is done with a Tomographic Gamma Scanner (TGS), Tomography can be used to extract 3D information about the density and gamma emissions of the item.
Gamma Cameras
Semiconductor detectors are used in some Gamma Cameras and Gamma imaging systems
| Physical sciences | Devices | Physics |
803808 | https://en.wikipedia.org/wiki/Irish%20elk | Irish elk | The Irish elk (Megaloceros giganteus), also called the giant deer or Irish deer, is an extinct species of deer in the genus Megaloceros and is one of the largest deer that ever lived. Its range extended across Eurasia during the Pleistocene, from Ireland (where it is known from abundant remains found in bogs) to Lake Baikal in Siberia. The most recent remains of the species have been radiocarbon dated to about 7,700 years ago in western Russia. Its antlers, which can span across are the largest known of any deer. It is not closely related to either living species called the elk, with it being widely agreed that its closest living relatives are fallow deer (Dama).
Taxonomy
Research history
The first scientific descriptions of the animal's remains were made by Irish physician Thomas Molyneux in 1695, who identified large antlers from Dardistown—which were apparently commonly unearthed in Ireland—as belonging to the elk (known as the moose in North America), concluding that it was once abundant on the island. It was first formally named as Alce gigantea by Johann Friedrich Blumenbach in his Handbuch der Naturgeschichte in 1799, with Alce being a variant of Alces, the Latin name for the elk. The original Blumenbach's description of Alce gigantea provides rather scant information about the species, specifying only that this particular kind of "fossil elk" comes from Ireland and is characterized by immense body size. According to Blumenbach, the distance between summits of giant deer antlers may attain . This particular feature mentioned by Blumenbach permitted to Roman Croitor to identify the type specimen of giant deer that was figured and described for the first time in Louthiana of Thomas Wright. The holotype of Megaloceros giganteus (Blumenbach, 1799) is a well-preserved male skull with exceptionally large antlers found in Dunleer environs (County Louth, Ireland). The type specimen of giant deer is currently exposed in Barmeath Castle where Thomas Wright first saw and described it.
French scientist Georges Cuvier documented in 1812 that the Irish elk did not belong to any species of mammal currently living, declaring it "le plus célèbre de tous les ruminans fossiles" (the most famous of all fossil ruminants). In 1827 Joshua Brookes, in a listing of his zoological collection, named the new genus Megaloceros (spelled Megalocerus in the earlier editions) in the following passage:
The etymology being from Greek: "great" + "horn, antler". The type and only species named in the description being Megaloceros antiquorum, based on Irish remains now considered to belong to M. giganteus, making the former a junior synonym. The original description was considered by Adrian Lister in 1987 to be inadequate for a taxonomic definition. In 1828 Brookes published an expanded list in the form of a catalogue for an upcoming auction, which included the Latin phrase "Cornibus deciduis palmatis" as a description of the remains. The 1828 publication was approved by International Commission on Zoological Nomenclature (ICZN) in 1977 as an available publication for the basis of zoological nomenclature. Adrian Lister in 1987 judged that "the phase "Cornibus deciduis palmatis" constitutes a definition sufficient under the [International Code of Zoological Nomenclature] (article 12) to validate Megalocerus." The original spelling of Megalocerus was never used after its original publication.In 1844 Richard Owen named another synonym of the Irish elk, including it within the newly named subgenus Megaceros, Cervus (Megaceros) hibernicus. This has been suggested to be derived from another junior synonym of the Irish elk described by J. Hart in 1825, Cervus megaceros. Despite being a junior synonym, Megaloceros remained in obscurity and Megaceros became the common genus name for the taxon. The combination "Megaceros giganteus" was in use by 1871. George Gaylord Simpson in 1945 revived the original Megaloceros name, which became progressively more widely used, until a taxonomic decision in 1989 by the ICZN confirmed the priority of Megaloceros over Megaceros, and Megaloceros to be the correct spelling.
Before the 20th century, the Irish elk, having evolved from smaller ancestors with smaller antlers, was taken as a prime example of orthogenesis (directed evolution), an evolutionary mechanism opposed to Darwinian evolution in which the successive species within the lineage become increasingly modified in a single undeviating direction, evolution proceeding in a straight line void of natural selection. Orthogenesis was claimed to have caused an evolutionary trajectory towards antlers that became larger and larger, eventually causing the species' extinction because the antlers grew to sizes which inhibited proper feeding habits and caused the animal to become trapped in tree branches. In the 1930s, orthogenesis was disputed by Darwinians led by Julian Huxley, who noted that antler size was not grossly large, and was proportional to body size. The currently favoured view is that sexual selection was the driving force behind the large antlers rather than orthogenesis or natural selection.
Evolution
M. giganteus belongs to the genus Megaloceros. Megaloceros has often been placed into the tribe Megacerini, alongside other genera often collectively referred to as "giant deer", like Sinomegaceros and Praemegaceros. The taxonomy of giant deer lacks consensus, with genus names used for species varying substantially between authors. The earliest possible record of the genus is a partial antler from the Early Pleistocene MN 17 (2.5–1.8 Ma) of Stavropol Krai in the North Caucasus of Russia, which were given the name of M. stavropolensis in 2016, however this species has been subsequently suggested to belong to Arvernoceros. or Sinomegaceros. The oldest generally accepted records of the genus are from the late Early Pleistocene. Other species often considered to belong to Megaloceros include the reindeer sized M. savini, which is known from early Middle Pleistocene (~700,000–450,000 years ago) localities in England, France, Spain and Germany, and the more recently described species M. novocarthaginiensis, which is known from late Early Pleistocene (0.9–0.8 Ma) localities in Spain, and the small M. matritensis endemic to the Iberian peninsula during the late Middle Pleistocene (~400,000 to 250,000 years ago), which overlaps chronologically with the earliest M. giganteus records. Jan van der Made proposed M. novocarthaginiensis , M. savini and M. matritensis to be sequential chronospecies, due to shared morphological characteristics not found in M. giganteus and gradual transition of morphological characters through time. M. savini and related species have also been suggested to comprise the separate genus Praedama by other authors. While the M. savini/Praedama lineage is often suggested to be closely related to M. giganteus, most authors agree that this group of deer is unlikely to be directly ancestral to M. giganteus.
The origin of M. giganteus remains unclear, and appears to lie outside Western Europe. Jan van der Made has suggested that remains of an indeterminate Megaloceros species from the late Early Pleistocene (~1.2 Ma) of Libakos in Greece are closer to M. giganteus than the M. novocarthaginiensis-savini-matritensis lineage due to the shared molarisation of the lower fourth premolar (P4). Croitor has suggested that M. giganteus is closely related to what was originally described as Dama clactoniana mugharensis (which he proposes be named Megaloceros mugharensis) from the Middle Pleistocene of Tabun Cave in Israel, due to similarities in the antlers, molars and premolars. The earliest possible records of M. giganteus comes from Homersfield, England thought to be about 450,000 years ago—though the dating is uncertain. The oldest securely dated Middle Pleistocene records are those from Hoxne, England, which have been dated to Marine Isotope Stage 11 (424,000 to 374,000 years ago), other Middle Pleistocene early records include Steinheim an der Murr, Germany, (classified as M. g. antecedens) about 400,000–300,000 years ago and Swanscombe, England. Most remains of the Irish elk are known from the Late Pleistocene. A large proportion of the known remains of M. giganteus are from Ireland, which mostly date to the Allerød oscillation near the end of the Late Pleistocene around 13,000 years ago. Over 100 individuals have been found in Ballybetagh Bog near Dublin.
Some authors have proposed that Late Pleistocene M. giganteus should be divided into several subspecies including M. giganteus ruffii and M. giganteus giganteus, based primarily on differences in antler morphology.
It has been historically thought that, because both have palmated antlers, the Irish elk and fallow deer (Dama spp.) are closely related, this is supported by several other morphological similarities, including the lack of upper canines, proportionally long braincase and nasal bones, and proportionally short front portion of the skull. In 2005, two fragments of mitochondrial DNA (mtDNA) from the cytochrome b gene were extracted and sequenced from 4 antlers and a bone, the mtDNA found that the Irish elk was nested within Cervus, and were inside the clade containing living red deer (Cervus elaphus). Based on this, the authors suggested that the Irish elk and red deer interbred. However, another study from the same year in the journal Nature utilising both fragmentary mitochondrial DNA and morphological data found that the Irish elk was indeed most closely related to Dama. The close relationship with Dama was supported by another cytochrome b study in 2006, a 2015 study involving the full mitochondrial genome, and by a 2017 morphological analysis of the bony labyrinth. The 2006 and 2017 studies also directly suggest that the results of the 2005 cytochrome b paper were the result of DNA contamination.
Cladogram of Cervidae based on mitochondrial DNA:
A study of mitochondrial genomes from Sinomegaceros from the Late Pleistocene of East Asia found that the mitochondrial genomes of Megaloceros giganteus were nested within those of Sinomegaceros, suggesting that the two lineages interbred after their initial split. Cladogram of Megaloceros and Sinomegaceros mitochondrial genomes following Xiao et al. 2023.
Description
The Irish elk stood about tall at the shoulders, and had large palmate (flat and broad) antlers, the largest of any known deer, with the largest specimens reaching over from tip to tip (though it is rare for specimens to exceed across) and in weight. The antlers are considerably larger than those of living moose, being on average over twice the volume of moose antlers. For body size, at about and up to or more, the Irish elk was the heaviest known cervine ("Old World deer"); and tied with the extant Alaska moose (Alces alces gigas) as the third largest known deer, after the extinct Cervalces latifrons and Cervalces scotti. The shape and span of the antlers varied significantly over time and space, likely reflecting some populations adaptation to forested environments. Compared to Alces, Irish elk appear to have had a more robust skeleton, with older and more mature Alces skeletons bearing some resemblance to those of prime Irish elk, and younger Irish elk resembling prime Alces. Likely due to different social structures, the Irish elk exhibits more marked sexual dimorphism than Alces, with Irish elk bucks being notably larger than does. In total, Irish elk bucks may have ranged from , with an average of , and does may have been relatively large, about 80% of buck size, or on average. The distinguishing characters of M. giganteus include concave frontals, proportionally long braincase, proportionally short front section of the skull (orbitofrontal region), alongside the absence of upper canines and the molarisation of the lower fourth premolar (P4). The skull and mandible of the Irish elk exhibit substantial thickening (pachyostosis), with the early and complete obliteration of cranial sutures.
Based on Upper Palaeolithic cave paintings, the Irish elk seems to have had overall light colouration, with a dark stripe running along the back, a stripe on either side from shoulder to haunch, a dark collar on the throat and a chinstrap, and a dark hump on the withers (between the shoulder blades). In 1989, American palaeontologist Dale Guthrie suggested that, like bison, the hump allowed a higher hinging action of the front legs to increase stride length while running. Valerius Geist suggested that the hump may have also been used to store fat. Localising fat rather than evenly distributing it may have prevented overheating while running or in rut during the summer.
Habitat
The Irish elk had a far-reaching range, extending from the Atlantic Ocean in the West to Lake Baikal in the East. Irish elk do not appear have extended northward onto the open mammoth steppe in Siberia, rather keeping to the boreal steppe-woodland environments, which consisted of scattered spruce and pine, as well as low-lying herbs and shrubs including grasses, sedges, Ephedra, Artemisia and Chenopodiaceae. The species appears to have had a degree of ecological plasticity, as during interglacial periods prior to the Holocene, the species was present in temperate forested environments in Europe. During these times, the species generally had less broad antlers than during glacial periods, likely as an adaptation to moving through forested environments.
Palaeobiology
Physiology
In 1998, Canadian biologist Valerius Geist hypothesised that the Irish elk was cursorial (adapted for running and stamina). He noted that the Irish elk physically resembled reindeer. The body proportions of the Irish elk are similar to those of the cursorial addax, oryx, and saiga antelope. These include the relatively short legs, the long front legs nearly as long as the hind legs, and a robust cylindrical body. Cursorial saiga, gnus, and reindeer have a top speed of over , and can maintain high speeds for up to 15 minutes.
Reproduction
At Ballybetagh Bog, over 100 Irish elk individuals were found, all small antlered bucks. This indicates that bucks and does segregated during at least winter and spring. Many modern deer species do this partly because males and females have different nutritional requirements and need to consume different types of plants. Segregation would also imply a polygynous society, with stags fighting for control over harems during rut. Because most of the individuals found were juvenile or geriatric and were likely suffering from malnutrition, they probably died from winterkill. Most Irish elk specimens known may have died from winterkill, and winterkill is the highest source of mortality among many modern deer species. Bucks generally suffer higher mortality rates because they eat little during the autumn rut. For rut, a lean stag normally may have fattened up to , and would burn through the extra fat over the next month.
Assuming a similar response to starvation as red deer, a large, healthy Irish elk stag with antlers would have had antlers under poor conditions; and an average sized Irish elk stag with antlers would have had antlers under poorer conditions, similar sizes to the moose. A similar change in a typical Irish elk population with prime stags having antlers would result in antler weights of or less in worsening climatic conditions. This is within the range of present-day wapiti/red deer (Cervus spp.) antler weights. Irish elk antlers vary widely in form depending upon the habitat, such as a compact, upright shape in closed forest environments. Irish elk likely shed their antlers and re-grew a new pair during mating season. Antlers generally require high amounts of calcium and phosphate, especially those for stags which have larger structures, and the massive antlers of Irish elk may have required much greater quantities. Stags typically meet these requirements in part from their bones, suffering from a condition similar to osteoporosis while the antlers are growing, and replenishing them from food plants after the antlers have grown in or reclaiming nutrients from shed antlers.
The large antlers have generally been explained as being used for male-male battle during mating season. They may have also been used for display, to attract females and assert dominance against rival males. A finite element analysis of the antlers suggested that during fighting, the antlers were likely to interlock around the middle tine, the high stress when interlocking on the distal tine suggests that the fighting was likely more constrained and predictable than among extant deer, likely involving twisting motions, as is known in extant deer with palmated antlers.
In deer, gestation time generally increases with body size. A doe may have had a gestation period of about 274 days. Based on this and patterns seen in modern deer, last year's antlers in Irish elk bucks were potentially shed in early March, peak antler growth in early June, completion by mid-July, shedding velvet (a layer of blood vessels on the antlers in-use while growing them) by late July, and the height of rut falling on the second week of August. Geist, believing the Irish elk to have been a cursorial animal, concluded that a doe would have to have produced nutrient-rich milk so that her calf would have enough energy and stamina to keep up with the herd.
Diet and life history
The mesodont (meaning neither high (hypsodont) or low (brachydont) crowned) condition of the teeth suggests that the species was a mixed feeder, being able to both browse and graze. Pollen remains from teeth found in the North Sea around 43,000 years old were found to be dominated by Artemisia and other Asteraceae, with minor Plantago, Helianthemum, Plumbaginaceae and willow (Salix). Another earlier specimen from the Netherlands (dating to the Eemian interglacial or early in the Last Glacial Period) was found to have pollen of Apiaceae, including cow parsley (Anthriscus sylvestris), cow parsnip/hogweed (Heracleum), water pennywort (Hydrocotyle), Asteraeceae, Filipendula, Symphytum and grass embedded with its teeth. A stable isotope analysis of the terminal Pleistocene Irish population suggests a grass and forb based diet, supplemented by browsing during stressed periods. Dental wear patterns of specimens from the late Middle and Late Pleistocene of Britain suggest a diet tending towards mixed feeding and grazing, but with a wide range including leaf browsing. Comparisons of δ15N between Irish elk and red deer at the Middle Pleistocene site of Schöningen in Germany suggest that grasses were a more important component of the former's diet relative to the latter.
Examination of histological sections of their long bones suggests that the species has relatively rapid growth rates, reaching skeletal maturity by around 6 years of age. Analysis of the cementum layers of their teeth suggests that Irish elk reached a maximum lifespan of at least 19 years, comparable to moose.
Based on the dietary requirements of red deer, a lean Irish elk stag would have needed to consume of fresh forage daily. Assuming antler growth occurred over a span of 120 days, a stag would have required 1,372 g (3 lb) of protein daily, as well as access to nutrient- and mineral-dense forage starting about a month before antlers began sprouting and continuing until they had fully grown. Such forage is not very common, and stags perhaps sought after aquatic plants in lakes. After antler growing, stags could probably satisfy their nutritional requirements in productive sedge lands bordered by willow and birch forests.
Gnaw marks on found on Irish elk bones indicates that they were preyed on or scavenged by cave hyenas.
Relationship with early humans
At a number of Middle Paleolithic sites, remains of M. giganteus have been found with cut marks indicating butchery by Neanderthals. These include Bolomor Cave in Spain, dating to around 180,000 years ago, and De Nadale Cave and Riparo del Broion in northern Italy, dating to 71-69,000 and 50-44,000 years ago, respectively. Other sites probably resulting from exploitation of Irish elk by Neanderthals include Abri du Maras in southeast France, dating to 55-40,000 years ago. A mandible from Ofatinţi, Moldova dating to either the Eemian or the early Late Pleistocene, has been noted for having "tool-made notches on its lateral side".
A handful of Irish elk depictions are known from the art of the Upper Paleolithic in Europe. However, these are much less abundant than the common red deer and reindeer depictions. Only a handful of examples of modern human interaction are known. Several M. giganteus bones from the Chatelperronian levels of the Labeko Koba site in Spain are noted for bearing puncture marks, which have been interpreted as anthropogenic. A terminal Pleistocene (13,710-13,215 cal BP) skull from Lüdersdorf, Germany is noted to have had the antler and facial part of the skull deliberately removed. A calcaneum from an associated lower hind limb from the early Holocene site of Sosnovy Tushamsky in Siberia is noted to have "two short and deep traces of cutting blows", which are interpreted as "clear evidence of butchery". The use of shed antler bases is also known, at the terminal Pleistocene (Allerød) Endingen VI site in Germany, a shed antler base appears to have been used in a way analogous to a lithic core to produce "blanks" for the manufacture of barbed projectile tips. A ring-like mark on a shed antler beam from the similarly aged Paderborn site in Germany has been suggested to be anthropogenic.
Extinction
Outside of the Irish Late Pleistocene, remains of Irish elk are uncommon, suggesting that they were usually rare in the areas where they did occur.
Historically, its extinction has been attributed to the encumbering size of the antlers, a "maladaptation" making fleeing through forests especially difficult for males while being chased by human hunters, or being too taxing nutritionally when the vegetation makeup shifted. In these scenarios, sexual selection by does for stags with large antlers would have contributed to decline.
However, antler size decreased through the Late Pleistocene and into the Holocene, and so may not have been the primary cause of extinction. A reduction in forest density in the Late Pleistocene and a lack of sufficient high-quality forage is associated with a decrease in body and antler size. Such resource constriction may have cut female fertility rates in half. Human hunting may have forced Irish elk into suboptimal feeding grounds.
The distribution of M. giganteus is thought to have been strongly controlled by climatic conditions. The range of the Irish elk appears to have collapsed during the Last Glacial Maximum (LGM), with few remains known between 27,500 and 14,600 years ago, and none between 23,300 and 17,500 years ago. Known remains substantially increase during the latest Pleistocene Bølling–Allerød Interstadial, where it appears to have re-colonized northern Europe, with abundant remains in the UK, Ireland, and Denmark, though its range contracted again during the following Younger Dryas, disappearing from northern Europe by the end of the period. A 2021 study found that M. giganteus saw a progressive decline in mitochondrial genome diversity beginning around 50,000 years ago, which accelerated during the LGM.
By the early Holocene, the range of the species had been dramatically reduced, with the youngest records in the eastern part of its range near Lake Baikal dating to around 10,700–10,400 years Before Present (BP), surviving latest in central part of its range within European Russia and Western Siberia. It is suggested that extinction was contributed to by further climatic changes transforming preferred open habitat into uninhabitable dense forest. The final demise may have been caused by several factors both on a continental and regional scale, including climate change and hunting. The youngest dates in this region from Kamyshlov in Western Siberia and Maloarkhangelsk, Oryol Oblast In European Russia date to around 7,700-7,600 years ago, and it is suggested that it likely became extinct shortly after this time. Lister and Stewart concluded in a study of the extinction of the Irish elk that "it seems clear that environmental factors, cumulatively over thousands of years, reduced giant deer populations to a highly vulnerable state. In this situation, even relatively low-level hunting by small human populations could have contributed to its extinction."
Modern significance
Due to the abundance of Irish elk remains in Ireland, a thriving trade in their bones existed there during the 19th century to supply museums and collectors. Skeletons and skulls with attached antlers were also prized ornaments in aristocratic homes. The remains of Irish elk were of high value: "In 1865, full skeletons might fetch £30, while particularly good heads with antlers could cost £15." with £15 being more than 30 weeks' wages for a low skilled worker at the time. Indeed Leeds Philosophical and Literary Society bought a full skeleton in 1847, from Glennon's in Dublin, for £38. This specimen, discovered at Lough Gur near Limerick, is still on display at Leeds City Museum.
| Biology and health sciences | Deer | Animals |
804032 | https://en.wikipedia.org/wiki/Monster%20truck | Monster truck | A monster truck is a specialized off-road vehicle with a heavy duty suspension, four-wheel steering, large-displacement V8 engines and oversized tires constructed for competition and entertainment uses. Originally created by modifying stock pickup trucks and sport utility vehicles (SUVs), they have evolved into purpose-built vehicles with tube-frame chassis and fiberglass bodies rather than metal. A competition monster truck is typically tall, and equipped with off-road tires.
Monster trucks developed in the late 1970s and came into the public eye in the early 1980s as side acts at popular motocross, tractor pulling, and mud bogging events, where they were used in car-crushing demonstrations. Today they are usually the main attraction with motocross, mud bogging, ATV racing, or demolition derbies as supporting events.
Events
Monster truck shows typically have two main segments; a race and a freestyle stunt driving competition, with an intermission at the midway point of the event.
Races are conducted as a single-elimination tournament on short, symmetrical tracks, which may include obstacles such as junk cars or dirt mounds. The length and complexity of the track can vary with the size of the venue, with courses in indoor arenas typically being shorter with fewer obstacles. This has drawn comparisons to professional drag racing, and depending on the promoter or size of the event, may even make use of a christmas tree starting procedure.
In freestyle events, each driver puts on a performance consisting of stunts such as obstacle jumps, backflips, wheelies, and doughnuts. A panel of judges assign points to each performance and the driver with the most points is declared the winner, though the exact scoring format can change from promoter to promoter.
Historically, additional vehicles for the drivers to crush, such as motor homes and school buses, were placed on the track specifically for the freestyle event, however incidents of debris flying into the stands and causing serious injuries have influenced most event promoters to turn away from such obstacles. Most freestyle courses now consist mostly of large mounds and ramps erected to allow the trucks to perform large jumps and wheelies upon landing. Freestyle performances have a set time limit and only one truck is allowed on the track at a time as a safety measure. Freestyle events are typically the final competition of a show, as damage to the trucks would make them unable to race.
History
In the late 1970s, heavily modified pickup trucks were becoming popular and the sports of mud bogging and truck pulling were gaining in popularity. Several truck owners had created lifted trucks to compete in such events, and soon competition to hold the title of "biggest truck" developed. The trucks which garnered the most national attention were Bob Chandler's Bigfoot, Everett Jasmer's USA-1, Fred Shafer and Jack Willman Sr.'s Bear Foot, and Jeff Dane's King Kong. At the time, the largest tires the trucks were running were in diameter.
In April 1981, Bob Chandler drove over junked cars in Bigfoot in what is often believed to be the first monster truck to crush cars. Chandler drove Bigfoot over a pair of cars in a field as a test of the truck's ability, and filmed it to use as a promotional tool in his four-wheel drive performance shop. An event promoter saw the video of the car crush and asked Chandler to do it in front of a crowd. Initially hesitant because of the "destructive" image that could be associated with Bigfoot, Chandler eventually agreed. After some smaller shows, Chandler performed the feat in the Pontiac Silverdome in 1982. At this show, Chandler also debuted a new version of Bigfoot with tires. At a prior event in the early 1980s, when Bigfoot was still running tires, Bob George, one of the owners of a motorsport promotion company named Truck-a-rama – later known as the United States Hot Rod Association (USHRA) – is said to have coined the phrase "monster truck" when referring to Bigfoot. The term "monster truck" became the generic name for all trucks with oversized tires.
King Kong and Bear Foot each followed Bigfoot to tires, and soon other monster trucks, such as King Krunch, Maddog, and Virginia Giant were being constructed. These early trucks were built on stock chassis which were heavily reinforced, used leaf spring suspension, a stock body, and heavy axles from military-specification vehicles to support the tires.
For most of the early 1980s, monster trucks performed primarily exhibitions as a side show to truck pulling or mud bogging events. In August 1983, Bigfoot and USA-1 competed in the first side-by-side monster truck race, which was filmed for the television show That's Incredible. By 1985 major promoters, such as the USHRA and TNT Motorsports, were racing monster trucks regularly. In 1988, TNT Motorsports created a series to establish the first national championship of monster truck racing; USA-1 and rookie driver Rod Litzau edged out Bigfoot, driven by Rich Hoosier, for the title.
In 1988, to standardize rules for truck construction and safety, Bob Chandler, Braden, and George Carpenter formed the Monster Truck Racing Association (MTRA). The MTRA created standard safety rules to govern monster trucks. The organization still plays a major role in the sport's development in the US and EU.
With racing taking precedence, several teams began to think in new ways as to how the trucks could be built. Towards the end of 1988, Gary Cook and David Morris debuted Equalizer, a truck with a combination of coil springs and shock absorbers as the main source of suspension rather than the standard of leaf springs and shock absorbers. In 1989, Jack Willman Sr., now with his own truck, Taurus, debuted a new truck which used a solid axle suspension system made of parallel four-link suspensions and coilovers that together weighed in at close to . However, the biggest innovation came from Chandler, also in 1989, when the CAD-designed Bigfoot #8 debuted featuring a full tubular chassis and a long-travel suspension system made of triangulated four-link suspensions, bump stops, limit straps, cantilevers, and shock absorbers charged with nitrogen gas. The truck revolutionized how monster trucks were built, and within a few years most top-level teams built similar vehicles.
In 1991, TNT was purchased by the USHRA and their points series were merged. The Special Events championship began to grow in popularity with teams as it had open qualifying spots which the invite-only USHRA championship did not have. The Special Events series lost its Pendaliner sponsorship in 1997. The short-lived ProMT series started in 2000.
Even though racing was dominant as a competition, USHRA events began having freestyle exhibitions as early as 1993. These exhibitions were developed as drivers, notably Dennis Anderson of the extremely popular Grave Digger, began asking for time to come out and perform if they lost in early rounds of racing. Promoters began to notice the popularity of freestyle among fans, and in 2000 the USHRA began holding freestyle as a judged competition at events, and now awards a freestyle championship.
Promoters
The USHRA's Monster Jam series, now owned by Feld Entertainment, is currently the largest, touring through the United States, Canada and select regions of Europe. Other promoters of monster truck events include the Toughest Monster Truck Tour, the Monster X Tour, and Hot Wheels Monster Trucks Live.
Truck construction
The first monster trucks built were pickup trucks and SUVs that were modified with larger suspension and larger tires. Today, trucks now have custom built tubular chassis, with four-link suspension to provide up to of clearance, and they also now have fiberglass bodies that attach to the chassis separately and are designed to be easily removed and easily replaced when damaged. The use of fiberglass panel bodies has allowed monster truck owners to develop a wide variety of thematic concept trucks that scarcely resemble the modified stock trucks that became early monster trucks.
Engines are now typically mounted behind the driver on most trucks and are typically supercharged, run on a methanol-based fuel, and have displacement of up to . Axles are mostly taken from either heavy-duty military trucks or road vehicles such as school buses, and are modified to have a planetary gear reduction at the hub to help turn the tires. All trucks have hydraulic steering in both the front and the rear (four wheel steering), with the front wheels controlled by the steering wheel and the rear wheels by a toggle switch.
The tires are typically "terra" tires used on farm equipment, and are of size in diameter, in width, and fit on diameter rims.
Most trucks use a modified or custom-designed automatic transmission, such as a Turbo 400, Powerglide, Ford C6 transmission, or a TorqueFlite 727. A limited number of trucks utilize a Lenco transmission, which traces its roots to drag racing. Most of the automatic transmissions are heavily modified with transbrakes, manual valve bodies, and heavy duty gear sets. Trucks running a Lenco use a centrifugal clutch as opposed to a torque converter, which are used in automatic transmissions. Lenco transmissions are usually found in two-speed or three speed configurations, and are commonly shifted using compressed carbon dioxide.
The trucks have many safety features, several of which are required just to run in the indoor arenas that the trucks frequent. Trucks are equipped with three shut-off switches: a remote ignition interrupt (RII), which allows event stewards to stop a truck remotely, a switch within the driver's reach in the cab, and another at the rear of the truck so that all electrical power may be shut off in the event of a rollover.
Many trucks are constructed with the driver sitting in the center of the cab for visibility. Most cabs are shielded with Lexan or comparable polycarbonate, which not only protects the driver from track debris, but also allows for increased visibility. Drivers are required to wear firesuits, safety harnesses, helmets, and head and neck restraints. Most moving parts on the truck are also shielded, and high pressure components have restraining straps, both in case of an explosion.
Accidents
Below is a list of accidents resulting in fatalities.
On January 16, 2009, at a Monster Jam event in Tacoma, Washington 6-year-old Sebastian Hizey was fatally injured when he was struck by flying debris from the truck Natural High. Hizey succumbed to his injuries the next morning.
On January 25, 2009, the monster truck Samson was involved in an accident in Madison, Wisconsin that caused the death of announcer George Eisenhart, Jr. after he accidentally stepped in front of Samson while it was moving.
On October 6, 2013, the monster truck Big Show plowed into a crowd of spectators in Chihuahua City, Mexico, killing 8 people and injuring 79 others in the Chihuahua monster truck accident. This is the deadliest monster truck incident in the history of the sport.
On September 28, 2014, a monster truck plowed into a crowd of spectators in the Haaksbergen monster truck accident, killing three people.
Guinness World Records
The world's biggest monster truck is Bigfoot 5, built in 1986, with tires that measure .
The world's longest monster truck is the Sin City Hustler, which measures long and was created by Brad and Jen Campbell in 2014.
The fastest speed record for a monster truck was achieved on August 6, 2022, by Joe Sylvester in Bad Habit at a speed of .
The longest ramp jump done by a monster truck was achieved in 2013 by Joey Sylvester in Bad Habit at a distance of .
The first monster truck backflip in a scored competition was achieved in 2010 by Cam McQueen in Nitro Circus.
The first monster truck front flip in a scored competition was achieved in 2017 by Lee O'Donnell in VP Racing Fuels' Mad Scientist at Monster Jam World Finals 18.
In June 2020, with touring suspended due to the COVID-19 pandemic, Monster Jam staged an event for pay-per-view called Monster Jam Breaking World Records in Bradenton, Florida. As part of the event, many new Guinness-recognized world records were set. These include the highest ramp jump by a monster truck by Krysten Anderson in Grave Digger at a height of , the most monster trucks jumped by a monster truck, by Adam Anderson in Megalodon, jumping over 8 trucks, the most donuts (spins) in a monster truck in one minute by Bari Musawwir in Zombie, spinning 44 times, the most consecutive donuts (spins), also by Musawwir, with 58, the longest stoppie (nose wheelie) by Tom Meents in Max-D, at , the farthest bicycle (side wheelie) by Ryan Anderson in Son Uva Digger, at , and the longest monster truck wheelie was by Adam Anderson in Grave Digger, at 190.46 m (624 ft 10.44 in).
| Technology | Motorized road transport | null |
808767 | https://en.wikipedia.org/wiki/Vacuum%20distillation | Vacuum distillation | Vacuum distillation or distillation under reduced pressure is a type of distillation performed under reduced pressure, which allows the purification of compounds not readily distilled at ambient pressures or simply to save time or energy. This technique separates compounds based on differences in their boiling points. This technique is used when the boiling point of the desired compound is difficult to achieve or will cause the compound to decompose. Reduced pressures decrease the boiling point of compounds. The reduction in boiling point can be calculated using a temperature-pressure nomograph using the Clausius–Clapeyron relation.
Laboratory-scale applications
Compounds with a boiling point lower than 150 °C typically are distilled at ambient pressure. For samples with high boiling points, short-path distillation apparatus is commonly employed. This technique is amply illustrated in Organic Synthesis.
Rotary evaporation
Rotary evaporation is a common technique used in laboratories to concentrate or isolate a compound from solution. Many solvents are volatile and can easily be evaporated using rotary evaporation. Even less volatile solvents can be removed by rotary evaporation under high vacuum and with heating. It is also used by environmental regulatory agencies for determining the amount of solvents in paints, coatings and inks.
Safety considerations
Safety is an important consideration when glassware is under vacuum pressure. Scratches and cracks can result in implosions when the vacuum is applied. Wrapping as much of the glassware with tape as is practical helps to prevent dangerous scattering of glass shards in the event of an implosion.
Industrial-scale applications
Industrial-scale vacuum distillation has several advantages. Close boiling mixtures may require many equilibrium stages to separate the key components. One tool to reduce the number of stages needed is to utilize vacuum distillation. Vacuum distillation columns (as depicted in Figures 2 and 3) typically used in oil refineries have diameters ranging up to about 14 meters (46 feet), heights ranging up to about 50 meters (164 feet), and feed rates ranging up to about 25,400 cubic meters per day (160,000 barrels per day).
Vacuum distillation can improve a separation by:
Prevention of product degradation or polymer formation because of reduced pressure leading to lower tower bottoms temperatures,
Reduction of product degradation or polymer formation because of reduced mean residence time especially in columns using packing rather than trays.
Increasing yield, and purity.
Vacuum distillation in petroleum refining
Petroleum crude oil is a complex mixture of hundreds of different hydrocarbon compounds generally having from 3 to 60 carbon atoms per molecule, although there may be small amounts of hydrocarbons outside that range. The refining of crude oil begins with distilling the incoming crude oil in a so-called atmospheric distillation column operating at pressures slightly above atmospheric pressure.
Vacuum distillation can also be referred to as "low-temperature distillation".
In distilling the crude oil, it is important not to subject the crude oil to temperatures above 370 to 380 °C because high molecular weight components in the crude oil will undergo thermal cracking and form petroleum coke at temperatures above that. Formation of coke would result in plugging the tubes in the furnace that heats the feed stream to the crude oil distillation column. Plugging would also occur in the piping from the furnace to the distillation column as well as in the column itself.
The constraint imposed by limiting the column inlet crude oil to a temperature of less than 370 to 380 °C yields a residual oil from the bottom of the atmospheric distillation column consisting entirely of hydrocarbons that boil above 370 to 380 °C.
To further distill the residual oil from the atmospheric distillation column, the distillation must be performed at absolute pressures as low as 10 to 40 mmHg / Torr (About 5% atmospheric pressure) so as to limit the operating temperature to less than 370 to 380 °C.
Figure 2 is a simplified process diagram of a petroleum refinery vacuum distillation column that depicts the internals of the column and Figure 3 is a photograph of a large vacuum distillation column in a petroleum refinery.
The 10 to 40 mmHg absolute pressure in a vacuum distillation column increases the volume of vapor formed per volume of liquid distilled. The result is that such columns have very large diameters.
Distillation columns such those in Images 1 and 2, may have diameters of 15 meters or more, heights ranging up to about 50 meters, and feed rates ranging up to about 25,400 cubic meters per day (160,000 barrels per day).
The vacuum distillation column internals must provide good vapor–liquid contacting while, at the same time, maintaining a very low-pressure increase from the top of the column top to the bottom. Therefore, the vacuum column uses distillation trays only where products are withdrawn from the side of the column (referred to as side draws). Most of the column uses packing material for the vapor–liquid contacting because such packing has a lower pressure drop than distillation trays. This packing material can be either structured sheet metal or randomly dumped packing such as Raschig rings.
The absolute pressure of 10 to 40 mmHg in the vacuum column is most often achieved by using multiple stages of steam jet ejectors.
Many industries, other than the petroleum refining industry, use vacuum distillation on a much smaller scale. Copenhagen-based Empirical Spirits, a distillery founded by former Noma chefs, uses the process to create uniquely flavoured spirits. Their flagship spirit, Helena, is created using Koji, alongside Pilsner Malt and Belgian Saison Yeast.
Large-scale water purification
Vacuum distillation is often used in large industrial plants as an efficient way to remove salt from ocean water, in order to produce fresh water. This is known as desalination. The ocean water is placed under a vacuum to lower its boiling point and has a heat source applied, allowing the fresh water to boil off and be condensed. The condensing of the water vapor prevents the water vapor from filling the vacuum chamber, and allows the effect to run continuously without a loss of vacuum pressure. The heat from condensation of the water vapor is removed by a heat sink, which uses the incoming ocean water as the coolant and thus preheats the feed of ocean water. Some forms of distillation do not use condensers, but instead compress the vapor mechanically with a pump. This acts as a heat pump, concentrating the heat from the vapor and allowing for the heat to be returned and reused by the incoming untreated water source. There are several forms of vacuum distillation of water, with the most common being multiple-effect distillation, vapor-compression desalination, and multi-stage flash distillation.
Molecular distillation
Molecular distillation is vacuum distillation below the pressure of 0.01 torr (1.3 Pa). 0.01 torr is one order of magnitude above high vacuum, where fluids are in the free molecular flow regime, i.e. the mean free path of molecules is comparable to the size of the equipment. The gaseous phase no longer exerts significant pressure on the substance to be evaporated, and consequently, the rate of evaporation no longer depends on pressure. That is, because the continuum assumptions of fluid dynamics no longer apply, mass transport is governed by molecular dynamics rather than fluid dynamics. Thus, a short path between the hot surface and the cold surface is necessary, typically by suspending a hot plate covered with a film of feed next to a cold plate with a line of sight in between.
Molecular distillation is used industrially for purification of oils.
Gallery
| Physical sciences | Phase separations | Chemistry |
809005 | https://en.wikipedia.org/wiki/Balloon%20%28aeronautics%29 | Balloon (aeronautics) | In aeronautics, a balloon is an unpowered aerostat, which remains aloft or floats due to its buoyancy. A balloon may be free, moving with the wind, or tethered to a fixed point. It is distinct from an airship, which is a powered aerostat that can propel itself through the air in a controlled manner.
Many balloons have a basket, gondola, or capsule suspended beneath the main envelope for carrying people or equipment (including cameras and telescopes, and flight-control mechanisms).
Aerostation
Aerostation is an obsolete term referring to ballooning and the construction, operation, and navigation of lighter-than-air vehicles. Tiberius Cavallo's The History and Practice of Aerostation was published in 1785. Other books were published on the subject including by Monck Mason. Dramatist Frederick Pilon wrote a play with aerostation as its title.
Principles
A balloon is conceptually the simplest of all flying machines. The balloon is a fabric envelope filled with a gas that is lighter than the surrounding atmosphere. As the entire balloon is less dense than its surroundings, it rises, taking along with it a basket, attached underneath, which carries passengers or payload. Although a balloon has no propulsion system, a degree of directional control is possible by making the balloon rise or sink in altitude to find favorable wind directions.
There are three main types of balloons:
The hot air balloon or Montgolfière obtains its buoyancy by heating the air inside the balloon; it has become the most common type.
The gas balloon or Charlière is inflated with a gas of lower molecular weight than the ambient atmosphere; most gas balloons operate with the internal pressure of the gas the same as the pressure of the surrounding atmosphere; a superpressure balloon can operate with the lifting gas at pressure that exceeds that of the surrounding air, with the objective of limiting or eliminating the loss of gas from day-time heating; gas balloons are filled with gases such as:
hydrogen – originally used extensively but, since the Hindenburg disaster, is now seldom used due to its high flammability;
coal gas – although giving around half the lift of hydrogen, extensively used during the nineteenth and early twentieth century, since it was cheaper than hydrogen and readily available;
helium – used today for all airships and most manned gas balloons;
other gases have included ammonia and methane, but these have poor lifting capacity and other safety defects and have never been widely used.
The Rozière type has both heated and unheated lifting gases in separate gasbags. This type of balloon is sometimes used for long-distance record flights, such as the recent circumnavigations, but is not otherwise in use.
Both the hot air, or Montgolfière, balloon and the gas balloon are still in common use. Montgolfière balloons are relatively inexpensive, as they do not require high-grade materials for their envelopes, and they are popular for balloonist sport activity.
Hot air balloons
The first balloon which carried passengers used hot air to obtain buoyancy and was built by the brothers Joseph and Etienne Montgolfier in Annonay, France in 1783: the first passenger flight was 19 September 1783, carrying a sheep, a duck, and a rooster.
The first tethered manned balloon flight was by a larger Montgolfier balloon, probably on 15 October 1783. The first free balloon flight was by the same Montgolfier balloon on 21 November 1783.
When heated, air expands, so a given volume of space contains less air. This makes it lighter and, if its lifting power is greater than the weight of the balloon containing it, it will lift the balloon upwards. A hot air balloon can only stay up while it has fuel for its burner, to keep the air hot enough.
The Montgolfiers' early hot air balloons used a solid-fuel brazier which proved less practical than the hydrogen balloons that had followed almost immediately, and hot air ballooning soon died out.
In the 1950s, the convenience and low cost of bottled gas burners led to a revival of hot air ballooning for sport and leisure.
The height or altitude of a hot air balloon is controlled by turning the burner up or down as needed, unlike a gas balloon where ballast weights are often carried so that they can be dropped if the balloon gets too low, and in order to land some lifting gas must be vented through a valve.
Gas balloons
A man-carrying balloon using the light gas hydrogen for buoyancy was made by Professor Jacques Charles and flown less than a month after the Montgolfier flight, on 1 December 1783. Gas balloons have greater lift for a given volume, so they do not need to be so large, and they can also stay up for much longer than hot air, so gas balloons dominated ballooning for the next 200 years. In the 19th century, it was common to use manufactured town gas (coal gas) to fill balloons; this was not as light as pure hydrogen gas, having about half the lifting power, but it was much cheaper and readily available.
Light gas balloons are predominant in scientific applications, as they are capable of reaching much higher altitudes for much longer periods of time. They are generally filled with helium. Although hydrogen has more lifting power, it is explosive in an atmosphere rich in oxygen. With a few exceptions, scientific balloon missions are unmanned.
There are two types of light-gas balloons: zero-pressure and superpressure. Zero-pressure balloons are the traditional form of light-gas balloon. They are partially inflated with the light gas before launch, with the gas pressure the same both inside and outside the balloon. As the zero-pressure balloon rises, its gas expands to maintain the zero pressure difference, and the balloon's envelope swells.
At night, the gas in a zero-pressure balloon cools and contracts, causing the balloon to sink. A zero-pressure balloon can only maintain altitude by releasing gas when it goes too high, where the expanding gas can threaten to rupture the envelope, or releasing ballast when it sinks too low. Loss of gas and ballast limits the endurance of zero-pressure balloons to a few days.
A superpressure balloon, in contrast, has a tough and inelastic envelope that is filled with light gas to pressure higher than that of the external atmosphere, and then sealed. The superpressure balloon cannot change size greatly, and so maintains a generally constant volume. The superpressure balloon maintains an altitude of constant density in the atmosphere, and can maintain flight until gas leakage gradually brings it down.
Superpressure balloons offer flight endurance of months, rather than days. In fact, in typical operation an Earth-based superpressure balloon mission is ended by a command from ground control to open the envelope, rather than by natural leakage of gas.
High-altitude balloons are used as high flying vessels to carry scientific instruments (like weather balloons), or reach near-space altitudes to take footage or photos of the earth. These balloons can fly over 100,000 feet (30.5 km) into the air, and are designed to burst at a set altitude where the parachute will deploy to safely carry the payload back to earth.
Cluster ballooning uses many smaller gas-filled balloons for flight.
Combination balloons
Early hot air balloons could not stay up for very long because they used a lot of fuel, while early hydrogen balloons were difficult to take higher or lower as desired because the aeronaut could only vent the gas or drop off ballast a limited number of times. Pilâtre de Rozier realised that for a long-distance flight such as crossing the English Channel, the aeronaut would need to make use of the differing wind directions at different altitudes. It would be essential therefore to have good control of altitude while still able to stay up for a long time. He developed a combination balloon having two gas bags, the Rozier balloon. The upper one held hydrogen and provided most of the steady lift. The lower one held hot air and could be quickly heated or cooled to provide the varying lift for good altitude control.
In 1785 Pilâtre de Rozier took off in an attempt to fly across the Channel, but shortly into the flight the hydrogen gas bag caught fire and de Rozier did not survive the ensuing accident. This earned de Rozier the title "The First to Fly and the First to Die".
It wasn't until the 1980s that technology was developed to allow safe operation of the Rozier type, for example by using non-flammable helium as the lifting gas, and several designs have successfully undertaken long-distance flights.
Tethering and kite balloons
As an alternative to free flight, a balloon may be tethered to allow reliable take off and landing at the same location. Some of the earliest balloon flights were tethered for safety, and since then balloons have been tethered for many purposes, including military observation and aerial barrage, meteorological and commercial uses.
The natural spherical shape of a balloon is unstable in high winds. Tethered balloons for use in windy conditions are often stabilised by aerodynamic shaping and connecting to the tether by a halter arrangement. These are called kite balloons.
A kite balloon is distinct from a kytoon, which obtains a portion of its lift aerodynamically.
History
Antecedents
Unmanned hot air balloons are mentioned in Chinese history. Zhuge Liang of the Shu Han kingdom, in the Three Kingdoms era (220–280 AD) used airborne lanterns for military signaling. These lanterns are known as Kongming lanterns ( ). The Mongolian army learned of the Kongming lantern from the Chinese and used it in Battle of Legnica during the Mongol invasion of Poland.
In 1709 the Brazilian-Portuguese cleric Bartolomeu de Gusmão made a balloon filled with heated air rise inside a room in Lisbon. On August 8, 1709, in Lisbon, Gusmão managed to lift a small balloon made of paper with hot air about four meters in front of king John V and the Portuguese court He also claimed to have built a balloon named (Big bird) and attempted to lift himself from Saint George Castle in Lisbon, landing about one kilometre away. However the claim of this feat remains uncertain, even though there is record of this flight in the source used by the FAI the exact distance and conditions of the flight are not confirmed.
The first modern balloons
Following Henry Cavendish's 1766 work on hydrogen, Joseph Black proposed that a balloon filled with hydrogen would be able to rise in the air.
The first recorded manned flight was made in a hot air balloon built by the Montgolfier brothers on 21 November 1783. The flight started in Paris and reached a height of 500 feet or so. The pilots, Jean-François Pilâtre de Rozier and François Laurent d'Arlandes, covered about in 25 minutes.
On 1 December 1783, Professor Jacques Charles and the Robert brothers made the first gas balloon flight, also from Paris. Their hydrogen-filled balloon flew to almost 2,000 feet (600 m), stayed aloft for over 2 hours and covered a distance of 27 miles (43 km), landing in the small town of Nesles-la-Vallée.
The first Italian balloon ascent was made by Count Paolo Andreani and two other passengers in a balloon designed and constructed by the three Gerli brothers, on 25 February 1784. A public demonstration occurred in Brugherio a few days later, on 13 March 1784, when the vehicle flew to a height of 1,537 metres (5,043 ft) and a distance of 8 kilometres (5.0 mi). On 28 March Andreani received a standing ovation at La Scala, and later a medal from Joseph II, Holy Roman Emperor.
De Rozier, together with Joseph Proust, took part in a further flight on 23 June 1784, in a modified version of the Montgolfiers' first balloon christened La Marie-Antoinette after the Queen. They took off in front of the King of France and King Gustav III of Sweden. The balloon flew north at an altitude of approximately 3,000 metres, above the clouds, travelling 52 km in 45 minutes before cold and turbulence forced them to descend past Luzarches, between Coye et Orry-la-Ville, near the Chantilly forest.
The first balloon ascent in Britain was made by James Tytler on 25 August 1784 at Edinburgh, Scotland, in a hot air balloon.
The first aircraft disaster occurred in May 1785 when the town of Tullamore, County Offaly, Ireland was seriously damaged when the crash of a balloon resulted in a fire that burned down about 100 houses, making the town home to the world's first aviation disaster. To this day, the town shield depicts a phoenix rising from the ashes.
Jean-Pierre Blanchard went on to make the first manned flight of a balloon in America on 9 January 1793, after touring Europe to set the record for the first balloon flight in countries including the Austrian Netherlands, Germany, the Netherlands and Poland. His hydrogen filled balloon took off from a prison yard in Philadelphia, Pennsylvania. The flight reached 5,800 feet (1,770 m) and landed in Gloucester County, New Jersey. President George Washington was among the guests observing the takeoff. Sophie Blanchard, married to Jean-Pierre, was the first woman to pilot her own balloon and the first woman to adopt ballooning as a career.
On 29 September 1804, Abraham Hopman became the first Dutchman to make a successful balloon flight in the Netherlands.
Gas balloons became the most common type from the 1790s until the 1960s. The French military observation balloon L'Intrépide of 1795 is the oldest preserved aircraft in Europe; it is on display in the Heeresgeschichtliches Museum in Vienna. Jules Verne wrote a short, non-fiction story, published in 1852, about being stranded aboard a hydrogen balloon.
The earliest successful balloon flight recorded in Australia was by William Dean in 1858. His balloon was gas-filled and travelled 30 km with two people aboard. On 5 January 1870, T. Gale, made an ascent from the Domain in Sydney. His balloon was 17 metres in length by 31 metres in circumference and his ascent, with him seated on the netting, took him about a mile before he landed in Glebe.
Henri Giffard also developed a tethered balloon for passengers in 1878 in the Tuileries Garden in Paris. The first tethered balloon in modern times was made in France at Chantilly Castle in 1994 by Aerophile SA.
Ballooning developed as a leisure activity. It was given a significant boost when Charles Green discovered that readily-available coal gas, then coming into urban use, gave half the lifting power of hydrogen, which had to be specially manufactured. In 1836 Green made an almost 500 mile long-distance flight from London, England to Weilberg in Germany.
Military use
The first military use of a balloon was at the Battle of Fleurus in 1794, when L'Entreprenant was used by the French Aerostatic Corps to watch the movements of the enemy. On 2 April 1794, an aeronauts corps was created in the French army; however, given the logistical problems linked with the production of hydrogen on the battlefield (it required constructing ovens and pouring water on white-hot iron), the corps was disbanded in 1799.
The first major use of balloons in the military occurred during the American Civil War with the Union Army Balloon Corps established in 1861.
During the Paraguayan War (1864–70), observation balloons were used by the Brazilian Army.
Balloons were used by the British Royal Engineers in 1885 for reconnaissance and observation purposes during the Bechuanaland Expedition and the Sudan Expedition. Although experiments in Britain had been conducted as early as 1863, a School of Ballooning was not established at Chatham, Medway, Kent until 1888. During the Anglo-Boer War (1899–1902), use was made of observation balloons. A balloon was kept inflated for 22 days and marched 165 miles into the Transvaal with the British forces.
Hydrogen-filled balloons were widely used during World War I (1914–1918) to detect enemy troop movements and to direct artillery fire. Observers phoned their reports to officers on the ground who then relayed the information to those who needed it. Balloons were frequently targets of opposing aircraft. Planes assigned to attack enemy balloons were often equipped with incendiary bullets, for the purpose of igniting the hydrogen.
The Aeronaut Badge was established by the United States Army in World War I to denote service members who were qualified balloon pilots. Observation balloons were retained well after the Great War, being used in the Russo-Finnish Wars, the Winter War of 1939–40, and the Continuation War of 1941–45.
During World War II the Japanese launched thousands of hydrogen "fire balloons" against the United States and Canada. In Operation Outward the British used balloons to carry incendiaries to Nazi Germany. During 2018, incendiary balloons and kites were launched from Gaza at Israel, burning some 12,000 dunams (3,000 acres) in Israel.
Large helium balloons are used by the South Korean government and private activists advocating freedom in North Korea. They float hundreds of kilometers across the border carrying news from the outside world, illegal radios, foreign currency and gifts of personal hygiene supplies. A North Korean military official has described it as "psychological warfare" and threatened to attack South Korea if their release continued.
Hot air returns
Ed Yost redesigned the hot air balloon in the late 1950s using rip-stop nylon fabrics and high-powered propane burners to create the modern hot air balloon. His first flight of such a balloon, lasting 25 minutes and covering 3 miles (5 km), occurred on 22 October 1960 in Bruning, Nebraska. Yost's improved design for hot air balloons triggered the modern sport balloon movement. Today, hot air balloons are much more common than gas balloons.
In the late 1970s the British hot air balloonist Julian Nott constructed a hot air balloon using technologies he believed would have been available to the Nazca culture of Peru some 1500 to 2000 years earlier, and demonstrated that it could fly. and again in 2003, Nott has speculated that the Nazca might have used it as a tool for designing the Nazca Lines. Nott also pioneered the use of hybrid energy, where solar power is a significant heat source, and in 1981 he crossed the English Channel.
Modern ballooning
In 2012, the Red Bull Stratos balloon took Felix Baumgartner to 128,100 ft. for a freefall jump from the stratosphere.
Sports
Commercial
Tethered gas balloons have been installed as amusement rides in Paris since 1999, in Berlin since 2000, in Disneyland Paris since 2005, in the San Diego Wild Animal Park since 2005, in Walt Disney World in Orlando since 2009, and the DHL Balloon in Singapore in 2006. Modern tethered gas balloons are made by Aerophile SAS.
Hot air balloons used in sport flying are sometimes made in special designs to advertise a company or product, such as the Chubb fire extinguisher illustrated.
Astronautics
Balloon satellites
A balloon in space uses internal gas pressure only to maintain its shape.
The Echo satellite was a balloon satellite launched into Earth orbit in 1960 and used for passive relay of radio communication. PAGEOS was launched in 1966 for worldwide satellite triangulation, allowing for greater precision in the calculation of different locations on the planet's surface.
Planetary probes
In 1984, the Soviet space probes Vega 1 and Vega 2 released two balloons with scientific experiments in the atmosphere of Venus. They transmitted signals for two days to Earth.
Ballooning records
On 19 October 1910, Alan Hawley and Augustus Post landed in the wilderness of Quebec, Canada after traveling for 48 hours and 1887.6 kilometers (1,173 mi) from St. Louis during the Gordon Bennett International Balloon Race, setting a distance record that held for more than 20 years. It took the men a week to hike out of the woods, during which time search parties had been mobilized and many had taken the pair for dead.
On 13 December 1913 through 17 December 1913 Hugo Kaulen stayed aloft for 87 hours. His record lasted until 1976.
On 27 May 1931, Auguste Piccard and Paul Kipfer became the first to reach the stratosphere in a balloon.
On 31 August 1933, Alexander Dahl took the first picture of the Earth's curvature in an open hydrogen gas balloon.
The helium-filled Explorer II balloon, piloted by US Army Air Corps officers Capt. Orvil A. Anderson, Maj. William E. Kepner and Capt. Albert W. Stevens, reached a new record height of 22,066 m (72,395 ft) on 11 November 1935. This followed the same crew's previous near-fatal plunge in July 1934 in a predecessor craft, Explorer, after its canopy ruptured just 190 m (624 ft) short (it transpired) of the then-current altitude record of 22,000 m (72,178 ft) set by the Soviet balloon Osoaviakhim-1.
In 1976, Ed Yost set 13 aviation world's records for distance traveled and amount of time aloft in his attempt to cross the Atlantic Ocean —solo— by balloon (3.938 km, 107:37 h).
In 1978, Ben Abruzzo and his team became the first to cross the Pacific Ocean in a hot air balloon.
The current absolute altitude record for manned balloon flight was set at 34,668 m (113,739 ft) on 4 May 1961 by Malcolm Ross and Victor Prather in the Strato-Lab V balloon payload launched from the deck of the USS Antietam in the Gulf of Mexico.
The previous record altitude for a manned balloon was set at 38,960.5 m (127,823 ft) by Felix Baumgartner in the Red Bull Stratos balloon launched from Roswell, New Mexico on Sunday, 14 October 2012.
The current record altitude for a manned balloon was set at 41,419.0 m (135,889.108 ft) by Alan Eustace on 24 October 2014 as part of the StratEx Space Dive project.
On 1 March 1999 Bertrand Piccard and Brian Jones set off in the balloon Breitling Orbiter 3 from Château d'Oex in Switzerland on the first non-stop balloon circumnavigation around the globe. They landed in Egypt after a 40,814 km (25,361 mi) flight lasting 19 days, 21 hours and 55 minutes.
The altitude record for an unmanned balloon is 53.0 kilometres (173,882 ft) in the mesosphere, reached with a volume of 60,000 cubic metres. The balloon was launched by JAXA on 25 May 2002 from Iwate Prefecture, Japan. This is the greatest height ever obtained by an atmospheric vehicle. Only rockets, rocket planes, and ballistic projectiles have flown higher.
In 2015, the two pilots Leonid Tiukhtyaev and Troy Bradley arrived safely in Baja California, Mexico after a journey of 10,711 km. The two men, originally from Russia and the United States of America respectively, started in Japan and flew with a helium balloon over the Pacific. In 160 hours, the balloon named "Two Eagles" arrived in Mexico, which is new distance and duration records for straight gas balloons.
| Technology | Aviation | null |
1886340 | https://en.wikipedia.org/wiki/Mixer%20%28appliance%29 | Mixer (appliance) | A mixer (also called a hand mixer or stand mixer depending on the type) is a kitchen device that uses a gear-driven mechanism to rotate a set of "beaters" in a bowl containing the food or liquids to be prepared by mixing them.
Mixers help automate the repetitive tasks of stirring, whisking or beating. When the beaters are replaced by a dough hook, a mixer may also be used to knead.
A mixer may be a handheld mechanism known as an eggbeater, a handheld motorized beater, or a drill mixer. Stand mixers vary in size from small counter top models for home use to large capacity commercial machines. Stand mixers create the mixing action by either rotating the mixing device vertically (planetary mixers), or by rotating the mixing container (spiral mixers).
History
The mixer with rotating parts was patented in 1856 by Baltimore, Maryland, tinner Ralph Collier. This was followed by E.P. Griffith's whisk patented in England in 1857. Another hand-turned rotary egg beater was patented by J.F. and E.P. Monroe in 1859 in the US. Their egg beater patent was one of the earliest bought up by the Dover Stamping Company, whose Dover egg beaters became a classic American brand. The term "Dover beater" was commonly in use in February 1929, as seen in this recipe from the Gazette newspaper of Cedar Rapids, IA, for "Hur-Mon Bavarian Cream," a whipped dessert recipe featuring gelatin, whipped cream, banana and gingerale. The Monroe design was also manufactured in England. In 1870, Turner Williams of Providence, R.I., invented another Dover egg beater model. In 1884, Willis Johnson of Cincinnati, Ohio, invented new improvements to the egg beater.
The first mixer with electric motor is thought to be the one invented by American Rufus Eastman in 1885. The Hobart Manufacturing Company was an early manufacturer of large commercial mixers, and they say a new model introduced in 1914 played a key role in the mixer part of their business. The Hobart KitchenAid and Sunbeam Mixmaster (first produced 1910) were two very early US brands of electric mixer. Domestic electric mixers were rarely used before the 1920s, when they were adopted more widely for home use.
In 1908 Herbert Johnston, an engineer for the Hobart Manufacturing Company, invented an electric standing mixer. His inspiration came from observing a baker mixing bread dough with a metal spoon; soon he was toying with a mechanical counterpart. By 1915, his 20-gallon (80 L) mixer was standard equipment for most large bakeries. In 1919, Hobart introduced the Kitchen Aid Food Preparer (stand mixer) for the home.
Older models of mixers originally listed each speed by name of operation (ex: Beat-Whip would be high speed if it is a 3-speed mixer); they are now listed by number.
Variants
Eggbeater
An eggbeater is a handheld device with a crank on the side geared to one or more beaters. The user grips the handle with one hand and operates the crank with the other, creating the rotary action.
Stand mixer
Stand mixers mount the motor driving the rotary action in a frame or stand which bears the weight of the device. Stand mixers are larger and have more powerful motors than their hand-held counterparts. They generally have a special bowl that is locked in place while the mixer is operating. A typical home stand mixer will include a wire whisk for whipping creams and egg whites; a flat beater for mixing batters; and a dough hook for kneading. Stand mixers are categorized as either spiral or planetary, based on whether or not the bowl is rotated.
Stand mixers are generally available in either counter top (also called bench) or floor models. Heavy duty commercial models can have bowl capacities well in excess of and weigh thousands of pounds (kilograms) but more typical home and light commercial models are equipped with bowls of around . Whether a mixer is a counter top or floor model depends on its size. Mixers that are in size or smaller tend to be counter top mixers, while larger mixers tend to be floor models due to their size and weight.
Spiral vs. planetary stand mixers
Spiral mixers are specialist tools for mixing dough. A spiral-shaped agitator counter-rotates while the powered bowl spins in the opposite direction. This method enables spiral mixers to mix the same size dough batch much quicker and with less under-mixed dough than a similarly powered planetary mixer. Spiral mixers can mix dough with less agitator friction than planetary mixers. This allows the dough to be mixed without increasing its temperature, ensuring the dough can rise properly. Spiral mixers are preferred for thicker products, such as dough for pizza, bagels or naan bread.
Planetary mixers consist of a bowl and an agitator. The bowl remains static, whilst the agitator is rapidly moved around the bowl to mix its contents. With the ability to mix a wide variety of ingredients, planetary mixers are more versatile than their spiral counterparts. Planetary mixers can be used to whip and blend, whereas spiral mixers cannot.
Hand mixer
A hand mixer is a hand-held mixing device. A handle is mounted over an enclosure containing the motor. The motor drives the beaters which are immersed in the food to perform the mixing action. The motor must be lightweight as it is supported by the user during use. The user may use any suitable kitchen container to hold the ingredients while mixing.
The first handheld electric mixer patent was submitted by Sunbeam Corporation in 1953 and granted in 1961.
Dough mixer
A dough mixer is used for household or industrial purposes. It is used for kneading large quantities of dough. It is electrical, having timers and various controls to suit the user's needs. Some features of dough blenders include high speed, low speed and bowl reverse (these can be combined into a programme) and a kneading bar in the centre of the bowl.
| Technology | Household appliances | null |
1886820 | https://en.wikipedia.org/wiki/Electronic%20component | Electronic component | An electronic component is any basic discrete electronic device or physical entity part of an electronic system used to affect electrons or their associated fields. Electronic components are mostly industrial products, available in a singular form and are not to be confused with electrical elements, which are conceptual abstractions representing idealized electronic components and elements. A datasheet for an electronic component is a technical document that provides detailed information about the component's specifications, characteristics, and performance. Discrete circuits are made of individual electronic components that only perform one function each as packaged, which are known as discrete components, although strictly the term discrete component refers to such a component with semiconductor material such as individual transistors.
Electronic components have a number of electrical terminals or leads. These leads connect to other electrical components, often over wire, to create an electronic circuit with a particular function (for example an amplifier, radio receiver, or oscillator). Basic electronic components may be packaged discretely, as arrays or networks of like components, or integrated inside of packages such as semiconductor integrated circuits, hybrid integrated circuits, or thick film devices. The following list of electronic components focuses on the discrete version of these components, treating such packages as components in their own right.
Classification
Components can be classified as passive, active, or electromechanic. The strict physics definition treats passive components as ones that cannot supply energy themselves, whereas a battery would be seen as an active component since it truly acts as a source of energy.
However, electronic engineers who perform circuit analysis use a more restrictive definition of passivity. When only concerned with the energy of signals, it is convenient to ignore the so-called DC circuit and pretend that the power supplying components such as transistors or integrated circuits is absent (as if each such component had its own battery built in), though it may in reality be supplied by the DC circuit. Then, the analysis only concerns the AC circuit, an abstraction that ignores DC voltages and currents (and the power associated with them) present in the real-life circuit. This fiction, for instance, lets us view an oscillator as "producing energy" even though in reality the oscillator consumes even more energy from a DC power supply, which we have chosen to ignore. Under that restriction, we define the terms as used in circuit analysis as:
Active components rely on a source of energy (usually from the DC circuit, which we have chosen to ignore) and usually can inject power into a circuit, though this is not part of the definition. Active components include amplifying components such as transistors, triode vacuum tubes (valves), and tunnel diodes.
Passive components cannot introduce net energy into the circuit. They also cannot rely on a source of power, except for what is available from the (AC) circuit they are connected to. As a consequence, they cannot amplify (increase the power of a signal), although they may increase a voltage or current (such as is done by a transformer or resonant circuit). Passive components include two-terminal components such as resistors, capacitors, inductors, and transformers.
Electromechanical components can carry out electrical operations by using moving parts or by using electrical connections.
Most passive components with more than two terminals can be described in terms of two-port parameters that satisfy the principle of reciprocity—though there are rare exceptions. In contrast, active components (with more than two terminals) generally lack that property.
Active components
Semiconductors
Transistors
Transistors were considered the invention of the twentieth century that changed electronic circuits forever. A transistor is a semiconductor device used to amplify and switch electronic signals and electrical power.
Field-effect transistors (FET)
MOSFET (metal–oxide–semiconductor FET) – by far the most widely manufactured electronic component (also known as MOS transistor)
PMOS (p-type MOS)
NMOS (n-type MOS)
CMOS (complementary MOS)
Power MOSFET
LDMOS (lateral diffused MOSFET)
MuGFET (multi-gate field-effect transistor)
FinFET (fin field-effect transistor)
TFT (thin-film transistor)
FeFET (ferroelectric field-effect transistor)
CNTFET (carbon nanotube field-effect transistor)
JFET (junction field-effect transistor) – N-channel or P-channel
SIT (static induction transistor)
MESFET (metal semiconductor FET)
HEMT (high-electron-mobility transistor)
Composite transistors
BiCMOS (bipolar CMOS)
IGBT (Insulated-gate bipolar transistor)
Other transistors
Bipolar junction transistor (BJT, or simply "transistor") – NPN or PNP
Photo transistor – amplified photodetector
Darlington transistor – NPN or PNP
Photo Darlington – amplified photodetector
Sziklai pair (compound transistor, complementary Darlington)
Tetrode transistor – is any transistor having four active terminals.
Thyristors
Silicon-controlled rectifier (SCR) – passes current only after triggered by a sufficient control voltage on its gate
TRIAC (TRIode for Alternating Current) – bidirectional SCR
Unijunction transistor (UJT)
Programmable Unijunction transistor (PUT)
SITh (static induction thyristor)
Diodes
Conduct electricity easily in one direction, among more specific behaviors.
Diode, rectifier, diode bridge
Schottky diode (hot carrier diode) – super fast diode with lower forward voltage drop
Zener diode – allows current to flow "backwards" when a specific set voltage is reached.
Transient voltage suppression diode (TVS), unipolar or bipolar – used to absorb high-voltage spikes
Varicap, tuning diode, varactor, variable capacitance diode – a diode whose AC capacitance varies according to the DC voltage applied.
Laser diode
Light-emitting diode (LED) – a diode that emits light
Photodiode – passes current in proportion to incident light
Avalanche photodiode – photodiode with internal gain
Solar Cell, photovoltaic cell, PV array or panel – produces power from light
DIAC (diode for alternating current), Trigger Diode, SIDAC) – often used to trigger an SCR
Constant-current diode
Step recovery diode
Tunnel diode - very fast diode based on quantum mechanical tunneling
Integrated circuits
Integrated Circuits can serve a variety of purposes, including acting as a timer, performing digital to analog conversion, performing amplification, or being used for logical operations.
Integrated circuit (IC)
MOS integrated circuit (MOS IC)
Hybrid integrated circuit (hybrid IC)
Mixed-signal integrated circuit
Three-dimensional integrated circuit (3D IC)
Digital electronics
Logic gate
Microcontroller
Analog circuit
Hall-effect sensor – senses a magnetic field
Current sensor – senses a current through it
Programmable devices
Programmable logic device
Field-programmable gate array (FPGA)
Complex programmable logic device (CPLD)
Field-programmable analog array (FPAA)
Optoelectronic devices
Opto-electronics
Opto-isolator, opto-coupler, photo-coupler – photodiode, BJT, JFET, SCR, TRIAC, zero-crossing TRIAC, open collector IC, CMOS IC, solid state relay (SSR)
Slotted optical switch, opto switch, optical switch
LED display – seven-segment display, sixteen-segment display, dot-matrix display
Display technologies
Current:
Filament lamp (indicator lamp)
Vacuum fluorescent display (VFD) (preformed characters, 7 segment, starburst)
Cathode-ray tube (CRT) (dot matrix scan, radial scan (e.g. radar), arbitrary scan (e.g. oscilloscope)) (monochrome & colour)
LCD (preformed characters, dot matrix) (passive, TFT) (monochrome, colour)
Neon (individual, 7 segment display)
LED (individual, 7 segment display, starburst display, dot matrix)
Split-flap display (numeric, preprinted messages)
Plasma display (dot matrix)
OLED (similar to an LCD, but each pixel generates its own light, can be made flexible or transparent)
Micro-LED (similar to OLED, but uses inorganic LEDs instead of organic ones, does not suffer from screen burn-in, however it cannot be made flexible or transparent)
Obsolete:
Incandescent filament 7 segment display (aka 'Numitron')
Nixie tube
Dekatron (aka glow transfer tube)
Magic eye tube indicator
Penetron (a 2 colour see-through CRT)
Vacuum tubes (valves)
A vacuum tube is based on current conduction through a vacuum (see Vacuum tube).
Diode or rectifier tube
Amplification
Triode
Tetrode
Pentode
Hexode
Pentagrid (Heptode)
Octode
Traveling-wave tube
Klystron
Oscillation
Magnetron
Reflex klystron (obsolete)
Carcinotron
Optical detectors or emitters
Phototube or photodiode – tube equivalent of semiconductor photodiode
Photomultiplier tube – phototube with internal gain
Cathode-ray tube (CRT) or television picture tube (obsolete)
Vacuum fluorescent display (VFD) – modern non-raster sort of small CRT display
Magic eye tube – small CRT display used as a tuning meter (obsolete)
X-ray tube – generates x-rays
Discharge devices
Gas discharge tube
Ignitron
Thyratron
Obsolete:
Mercury arc rectifier
Voltage regulator tube
Nixie tube
Power sources
Sources of electrical power:
Battery – acid- or alkali-based power supply.
Fuel cell – an electrochemical generator
Power supply – usually a main hook-up
Photovoltaic device – generates electricity from light
Thermoelectric generator – generates electricity from temperature gradients
Electrical generator – an electromechanical power source
Piezoelectric generator - generates electricity from mechanical strain
Van de Graaff generator - generates electricity from friction
Passive components
Components incapable of controlling current by means of another electrical signal are called passive devices. Resistors, capacitors, inductors, and transformers are all considered passive devices.
Resistors
Pass current in proportion to voltage (Ohm's law) and oppose current.
Resistor – fixed value
Power resistor – larger to safely dissipate heat generated
SIP or DIP resistor network – array of resistors in one package
Variable resistor
Rheostat – two-terminal variable resistor (often for high power)
Potentiometer – three-terminal variable resistor (variable voltage divider)
Trim pot – small potentiometer, usually for internal adjustments
Thermistor – thermally sensitive resistor whose prime function is to exhibit a large, predictable and precise change in electrical resistance when subjected to a corresponding change in body temperature.
Humistor – humidity-varied resistor
Photoresistor
Memristor
Varistor, Voltage-dependent resistor, MOV – Passes current when excessive voltage is present
Resistance wire, Nichrome wire – wire of high-resistance material, often used as a heating element
Heater – heating element
Capacitors
Capacitors store and release electrical charge. They are used for filtering power supply lines, tuning resonant circuits, and for blocking DC voltages while passing AC signals, among numerous other uses.
Capacitor
Integrated capacitors
MIS capacitor
Trench capacitor
Fixed capacitors
Ceramic capacitor
Film capacitor
Electrolytic capacitor
Aluminum electrolytic capacitor
Tantalum electrolytic capacitor
Niobium electrolytic capacitor (Columbium capacitor)
Polymer capacitor, OS-CON
Supercapacitor (Electric double-layer capacitor)
Nanoionic supercapacitor
Lithium-ion capacitor
Mica capacitor
Vacuum capacitor
Variable capacitor – adjustable capacitance
Tuning capacitor – variable capacitor for tuning a radio, oscillator, or tuned circuit
Trimmer capacitor – small variable capacitor for seldom or rare adjustments of LC-circuits
Vacuum variable capacitor
Capacitors for special applications
Power capacitor
Safety capacitor
Filter capacitor
Light-emitting capacitor (LEC)
Motor capacitor
Photoflash capacitor
Reservoir capacitor / Bulk capacitor
Coupling capacitor
Decoupling capacitor / Buffer capacitor
Bypass capacitor
Pull capacitor / Padding capacitor
Backup capacitor
Switched capacitor
Feedthrough capacitor
Capacitor network (array)
Varicap diode – AC capacitance varies according to the DC voltage applied
Integrated passive devices
Integrated passive devices are passive devices integrated within one distinct package. They take up less space than equivalent combinations of discrete components.
Magnetic (inductive) devices
Electrical components that use magnetism in the storage and release of electrical charge through current:
Inductor, coil, choke
Variable inductor
Saturable inductor
Transformer
Magnetic amplifier (toroid)
ferrite impedances, beads
Motor / Generator
Solenoid
Loudspeaker and microphone
Memristor
Electrical components that pass charge in proportion to magnetism or magnetic flux, and have the ability to retain a previous resistive state, hence the name of Memory plus Resistor.
Memristor
Networks
Components that use more than one type of passive component:
RC network – forms an RC circuit, used in snubbers
LC Network – forms an LC circuit, used in tunable transformers and RFI filters.
Transducers, sensors, detectors
Transducers generate physical effects when driven by an electrical signal, or vice versa.
Sensors (detectors) are transducers that react to environmental conditions by changing their electrical properties or generating an electrical signal.
The transducers listed here are single electronic components (as opposed to complete assemblies), and are passive (see Semiconductors and Tubes for active ones). Only the most common ones are listed here.
Audio
Loudspeaker – Electromagnetic or piezoelectric device to generate full audio
Buzzer – Electromagnetic or piezoelectric sounder to generate tones
Position, motion
Linear variable differential transformer (LVDT) – Magnetic – detects linear position
Rotary encoder, Shaft Encoder – Optical, magnetic, resistive or switches – detects absolute or relative angle or rotational speed
Inclinometer – Capacitive – detects angle with respect to gravity
Motion sensor, Vibration sensor
Flow meter – detects flow in liquid or gas
Force, torque
Strain gauge – Piezoelectric or resistive – detects squeezing, stretching, twisting
Accelerometer – Piezoelectric – detects acceleration, gravity
Thermal
Thermocouple, thermopile – Wires that generate a voltage proportional to delta temperature
Thermistor – Resistor whose resistance changes with temperature, up PTC or down NTC
Resistance Temperature Detector (RTD) – Wire whose resistance changes with temperature
Bolometer – Device for measuring the power of incident electromagnetic radiation
Thermal cutoff – Switch that is opened or closed when a set temperature is exceeded
Magnetic field (see also Hall Effect in semiconductors)
Magnetometer, Gauss meter
Humidity
Hygrometer
Electromagnetic, light
Photo resistor – Light dependent resistor (LDR)
Antennas
Antennas transmit or receive radio waves
Elemental dipole
Yagi
Phased array
Loop antenna
Parabolic dish
Log-periodic dipole array
Biconical
Feedhorn
Assemblies, modules
Multiple electronic components assembled in a device that is in itself used as a component
Oscillator
Display devices
Liquid crystal display (LCD)
Digital voltmeters
Filter
Prototyping aids
Wire-wrap
Breadboard
Electromechanical devices
Piezoelectric devices, crystals, resonators
Passive components that use piezoelectric effect:
Components that use the effect to generate or filter high frequencies
Crystal – a ceramic crystal used to generate precise frequencies (See the Modules class below for complete oscillators)
Ceramic resonator – Is a ceramic crystal used to generate semi-precise frequencies
Ceramic filter – Is a ceramic crystal used to filter a band of frequencies such as in radio receivers
surface acoustic wave (SAW) filters
Components that use the effect as mechanical transducers.
Ultrasonic motor – Electric motor that uses the piezoelectric effects
For piezo buzzers and microphones, see the Transducer class below
Microelectromechanical systems
Microelectromechanical systems
Accelerometer
Digital micromirror device
Terminals and connectors
Devices to make electrical connection
Terminal
Connector
Socket
Screw terminal, Terminal Blocks
Pin header
Cable assemblies
Electrical cables with connectors or terminals at their ends
Power cord
Patch cord
Test lead
Switches
Components that can pass current ("closed") or break the current ("open"):
Switch – Manually operated switch
Electrical description: SPST, SPDT, DPST, DPDT, NPNT (general)
Technology: slide switches, toggle switches, rocker switches, rotary switches, pushbutton switches
Keypad – Array of pushbutton switches
DIP switch – Small array of switches for internal configuration settings
Footswitch – Foot-operated switch
Knife switch – Switch with unenclosed conductors
Micro switch – Mechanically activated switch with snap action
Limit switch – Mechanically activated switch to sense limit of motion
Mercury switch – Switch sensing tilt
Centrifugal switch – Switch sensing centrifugal force due to rate of rotation
Relay or contactor – Electro-mechanically operated switch (see also solid state relay above)
Reed switch – Magnetically activated switch
Thermostat – Thermally activated switch
Humidistat – Humidity activated switch
Circuit breaker – Switch opened in response to excessive current: a resettable fuse
Disconnector – Switch used in high- and medium-voltage applications for maintenance of other devices or isolation of circuits
Transfer switch – Switch that toggles a load between two sources
Protection devices
Passive components that protect circuits from excessive currents or voltages:
Fuse – over-current protection, one time use
Circuit breaker – resettable fuse in the form of a mechanical switch
Resettable fuse or PolySwitch – circuit breaker action using solid state device
Ground-fault protection or residual-current device – circuit breaker sensitive to mains currents passing to ground
Metal oxide varistor (MOV), surge absorber, TVS – Over-voltage protection
Inrush current limiter – protection against initial Inrush current
Gas discharge tube – protection against high voltage surges
Spark gap – electrodes with a gap to arc over at a high voltage
Lightning arrester – spark gap used to protect against lightning strikes
Recloser – automatic switch that opens on an overcurrent (fault) condition, then closes to check if the fault is cleared, and repeats this process a specified number of times before maintaining the open position until it is manually closed
Arc-fault circuit interrupter – circuit breaker that protects against arcs
Network protector – protective device that disconnects a distribution transformer when energy flow reverses direction
Magnetic starter – electromechanical switch used in motors
Mechanical accessories
Enclosure (electrical)
Heat sink
Fan
Other
Printed circuit boards
Lamp
Waveguide
Obsolete
Carbon amplifier (see Carbon microphones used as amplifiers)
Carbon arc (negative resistance device)
Dynamo (historic rf generator)
Coherer
Standard symbols
On a circuit diagram, electronic devices are represented by conventional symbols. Reference designators are applied to the symbols to identify the components.
| Technology | Components | null |
1886997 | https://en.wikipedia.org/wiki/Stearyl%20alcohol | Stearyl alcohol | Stearyl alcohol, or 1-octadecanol, is an organic compound classified as a saturated fatty alcohol with the formula CH3(CH2)16CH2OH. It takes the form of white granules or flakes, which are insoluble in water. It has a wide range of uses as an ingredient in lubricants, resins, perfumes, and cosmetics. It is used as an emollient, emulsifier, and thickener in ointments, and is widely used as a hair coating in shampoos and hair conditioners. Stearyl heptanoate, the ester of stearyl alcohol and heptanoic acid (enanthic acid), is found in most cosmetic eyeliners. Stearyl alcohol has also found application as an evaporation suppressing monolayer when applied to the surface of water.
Stearyl alcohol is prepared from stearic acid or certain fats by the process of catalytic hydrogenation. It has low toxicity.
| Physical sciences | Alcohols | Chemistry |
1887046 | https://en.wikipedia.org/wiki/Coreidae | Coreidae | Coreidae is a large family of predominantly sap-sucking insects in the Hemipteran suborder Heteroptera. The name "Coreidae" derives from the genus Coreus, which derives from the Ancient Greek () meaning bedbug.
As a family, the Coreidae are cosmopolitan, but most of the species are tropical or subtropical.
Common names and significance
The common names of the Coreidae vary regionally. Leaf-footed bug refers to leaf-like expansions on the legs of some species, generally on the hind tibiae. In North America, the pest status of species such as Anasa tristis on squash plants and other cucurbits gave rise to the name squash bugs. The Coreidae are called twig-wilters or tip-wilters in parts of Africa and Australia because many species feed on young twigs, injecting enzymes that macerate the tissues of the growing tips and cause them to wilt abruptly.
Morphology and appearance
The Coreidae commonly are oval-shaped, with antennae composed of four segments, numerous veins in the membrane of the fore wings, and externally visible repugnatorial stink glands. They vary in size from 7 to 45 mm long, which implies that the family includes some of the biggest species of Heteroptera. The body shape is quite variable; some species are broadly oval, others are elongated with parallel sides, and a few are slender. Many species with the "leaf-footed" tibiae are very slender with conspicuous expansions of the hind tibiae, but some robust species also have decided expansions. Some species are covered with spines and tubercles. As an example of these, the tribe Phyllomorphini Mulsant & Rey, 1870, are strikingly aberrant, with thin legs, spiny bristles, and laciniate outlines and adornments.
Many of the more robust species have grossly enlarged, thickened, and bowed hind femora armed with spikes on the inner edge, and with hind tibiae to match, though the enlargement of the tibiae is less exaggerated.
In the nymphs, the openings of the two repugnatorial stink glands of the Coreidae are visible as two projections or spots on the medial line of the dorsal surface of the abdomen, one at the anterior and one at the posterior edge of the fifth abdominal tergite above the glands inside. During the final ecdysis, the anatomy is rearranged and the glands end up in the metathorax, opening laterally through ostioles between the mesothoracic and metathoracic pleura.
Biology and habits
The Coreidae generally feed on the sap of plants. Some species reportedly are actively carnivorous, but material evidence is lacking, and in the field, some are easy to confuse with some species of the Reduviidae, so doubt has been cast on the veracity of the claims.
Some Coreidae, such as Phyllomorpha laciniata, exhibit parental care by carrying their eggs. This behaviour significantly improves the eggs' chances of avoiding the attacks of parasitoids.
Taxonomy and systematics
The Coreidae are placed in the order Hemiptera and closely related to the families Alydidae, Hyocephalidae, Rhopalidae, and Stenocephalidae. Together, these five families form the superfamily Coreoidea. The family is large, with more than 1,900 species in over 270 genera.
Most taxonomists dealing with the Coreidae divide the family into three or four subfamilies. Numerous tribes of the Coreinae have previously been proposed for elevation to subfamily rank, for example, the Agriopocorini, Colpurini, Phyllomorphini, and Procamptini, but the only one of these changes that at least a significant minority of researchers accepted is the elevation of the Agriopocorinae, and more recent reviews tend to treat them as a tribe again, recognizing only the three subfamilies known by 1867 plus Hydarinae.
The family has been demonstrated to be non-monophyletic, as Hydarinae and Pseudophloeinae are more closely related to Alydidae than to other coreids.
Selected genera
The following genera are included in the family Coreidae:
Coreinae Leach, 1815
Acanthocephala Laporte, 1833
Acanthocerus Palisot, 1818
Acanthocoris Amyot & Serville, 1843
Agathyrna Stål, 1861
Althos Kirkaldy, 1904
Anasa Amyot & Serville, 1843
Anisoscelis Latreille, 1829
Anoplocnemis Stål, 1873
Aurelianus Distant, 1902
Brachytes Westwood, 1842
Canungrantmictis Brailovsky, 2002
Carlisis Stål, 1858
Catorhintha Stål, 1859
Centrocoris Kolenati, 1845
Ceratopachys Westwood, 1842
Cercinthus Stål, 1860
Chariesterus Laporte, 1833
Chelinidea Uhler, 1863
Cimolus Stål, 1862
Coreus Leach, 1815
Dalader Amyot & Serville, 1843
Elasmopoda Stål, 1873
Enoplops Amyot & Serville, 1843
Eubule Stål, 1868
Euthochtha Mayr, 1865
Fracastorius Distant, 1902
Gelonus Stål, 1866
Gonocerus Berthold, 1827
Helcomeria Stål, 1873
Himella Dallas, 1852
Holhymenia Lepeletier & Serville, 1825
Homoeocerus Burmeister, 1835
Hygia Uhler, 1861
Hypselonotus Hahn, 1833
Leptoglossus Guérin-Méneville, 1831 – conifer seed bugs
Mamurius Stål, 1862
Menenotus Laporte, 1832
Mictis Leach, 1814
Molipteryx Kiritshenko, 1916
Mozena Amyot & Serville, 1843
Namacus Amyot & Serville, 1843
Narnia Stål, 1862
Neaira Linnavuori, 1973
Nematopus Berthold, 1827
Nisoscolopocerus Barber, 1928
Oannes Distant, 1911
Pachylis Le Peletier & Serville, 1825
Phyllomorpha Laporte 1833
Physomerus Burmeister, 1835
Piezogaster Amyot & Serville, 1843
Plectropoda Bergroth, 1894
Pomponatius Distant, 1904
Prionolomia Stål, 1873
Pseudotheraptus Brown, 1955
Sagotylus Mayr, 1865
Savius Stål, 1862
Scolopocerus Uhler, 1875
Sephina Amyot & Serville, 1843
Sethenira Spinola, 1837
Spartocera Laporte, 1833
Spathocera Stein, 1860
Syromastus Berthold, 1827
Thasus Stål, 1865
Vazquezitocoris Brailovsky, 1990
Villasitocoris Brailovsky, 1990
Wolfius Distant, 1902
Zicca Amyot & Serville, 1843
Hydarinae Stål, 1873
Madura Stål, 1860
Meropachyinae Stål, 1867
Merocoris Perty, 1833
Pseudophloeinae Stål, 1867
Arenocoris Hahn, 1834
Bathysolen Fieber, 1860
Bothrostethus Fieber 1860
Ceraleptus Costa, 1847
Clavigralla Spinola, 1837
Coriomeris Westwood, 1842
Nemocoris Sahlberg, 1848
Ulmicola Kirkaldy, 1909
Gallery
| Biology and health sciences | Hemiptera (true bugs) | Animals |
1887197 | https://en.wikipedia.org/wiki/Pelican%20eel | Pelican eel | The pelican eel (Eurypharynx pelecanoides) is a deep-sea eel. It is the only known member of the genus Eurypharynx and the family Eurypharyngidae. It belongs to the "saccopharyngiforms", members of which were historically placed in their own order, but are now considered true eels in the order Anguilliformes. The pelican eel has been described by many synonyms, yet nobody has been able to demonstrate that more than one species of pelican eel exists. It is also referred to as the gulper eel (which can also refer to members of the related genus Saccopharynx), pelican gulper, and umbrella-mouth gulper. The specific epithet pelecanoides refers to the pelican, as the fish's large mouth is reminiscent of that of the pelican.
Description
The morphology of pelican eel specimens can be hard to describe because they are so fragile that they become damaged when they are recovered from the deep sea's immense pressure. However, certain observations about the physical characteristics have been noted from studied specimen.
The pelican eel's most notable feature is its large mouth, which is much larger than its body. The mouth is loosely hinged, and can be opened wide enough to swallow a fish much larger than the eel itself. The lower jaw is hinged at the base of the head, with no body mass behind it, making the head look disproportionately large. Its jaw is so large that it is estimated to be about a quarter of the total length of the eel itself.
While typically in a folded state, the pelican eel's mouth has the capacity to change to an inflated shape when hunting, giving the mouth its notably massive appearance. This transformation is possible due to the dual-mode biological morphing mechanism that takes place: geometric unfolding of the mouth followed by stretching. When the pelican eel is in pursuit of prey and opens its mouth, the head and jaw structure unfold and spread horizontally, This un-spreading event is followed by the inflation of the mouth. The inflation is made possible given the highly stretchable skin of the head, an additional characteristic that enables the eel to partake in this mechanism and engage in lunge feeding to consume large amounts of prey. When it feeds on prey, water that is ingested is expelled via the gills.
Pelican eels are smaller-sized eels. They grow to about in length, though lengths of are plausible. Like most eels, E. pelecanoides lacks pelvic fins and scales. Otherwise, the pelican eel is very different in appearance from typical eels. Instead of having a swim bladder, the pelican eel has an aglomerular kidney that is thought to have a role in maintaining the gelatinous substance filling the "lymphatic spaces" that are found around the vertebrae. It has been hypothesized that these gelatinous substance filled "lymphatic spaces" could function in a similar way to a swim bladder. Furthermore, the muscle segment shape of the pelican eel is different. Its muscle segments have a "V-shape", while other fish have "W-shaped" muscle segments. Pelican eels are also unusual because the ampullae of the lateral line system project from the body, rather than being contained in a narrow groove; this may increase its sensitivity.
Unlike many other deep sea creatures, the pelican eel has very small eyes. For reference, the horizontal eye size diameter of a male pelican eel specimen was measured to be . It is believed that the eyes evolved to detect faint traces of light rather than form images.
The pelican eel also has a very long, whip-like tail that it uses for movement and for communication via bioluminescence. Specimens that have been brought to the surface in fishing nets have been known to have their long tails tied into several knots. The end of the tail bears a complex organ with numerous tentacles, which glows pink and gives off occasional bright-red flashes. The colors on its tail are displayed through its light-emitting photophores. This is presumably a lure to attract prey, although its presence at the far end of the body from the mouth suggests the eel may have to adopt an unusual posture to use it effectively.
Pelican eels are black or olive and some subspecies may have a thin lateral white stripe. The coloration of E. pelecanoides is especially dark because this species exhibits ultra-black camouflage. This special pigmentation, which reflects less than 0.5% of light, allows these eels to be cloaked in darkness in their low light environments. Ultra-black camouflage allows these bathypelagic eels to evade predators and hide from prey.
Pelican eels display sexual dimorphism with the largest morphological difference in the structure of the nasal rosette. In female pelican eels, the nasal rosette is hardly noticeable whereas male pelican eels exhibit a larger nasal rosette. The male's nasal rosette is bulb-shaped and contains larger anterior and posterior nostrils. Sexual dimorphism is thought to aid with locating a potential mate in the bathypelagic zone.
Feeding habits and diet
Pelican eels have developed adaptations and feeding patterns to help them survive in their low biomass environment. Recent studies have shown that pelican eels are active participants in their pursuit of food, rather than passively waiting for prey to fall into their large mouths. They are hypothesized to exhibit lunge-feeding through the expansion of their mandible and upper jaw. Furthermore, their stomach can stretch and expand to accommodate large meals, although analysis of stomach contents suggests they primarily eat small crustaceans. Despite the great size of the jaws, which occupy about a quarter of the animal's total length, it has only tiny teeth, which would not be consistent with a regular diet of large fish.
The large mouth may be an adaptation to allow the eel to eat a wider variety of prey when food is scarce. The eel can swim into large groups of shrimp or other crustaceans with its mouth closed, opening wide as it closely approaches prey, scooping them up to be swallowed. The pelican eel is also known to feed on cephalopods (squid) and other small invertebrates. When the eel gulps its prey into its massive jaws, it also takes in a large amount of water, which is then slowly expelled through its gill slits. Pelican eels themselves are preyed upon by lancetfish and other deep sea predators. The pelican eel is not known to undergo vertical diurnal migration like other eels.
Observations of gut contents and teeth morphology indicate that Eurypharynx pelecanoides larva, categorized as a type of leptocephali, feed on marine snow. Organisms, such as thraustochytrids and hydrozoan tissue, were consumed by these larva in a grouped manner such as they would be found in marine snow. Furthermore, the lesser number, larger size, and inwardly-pointing direction of leptocephali larval teeth point indicate that pelican eel larva rely on marine snow as a source of nutrients. As leptocephali develop into their mature form these distinct teeth were replaced by more, smaller teeth. This particular observation may explain a shift in the size of leptocephali heads, such as E. pelecanoides, in comparison to their food source as they mature.
Reproduction and life cycle
Not much is known about the reproductive habits of the pelican eel. Similar to other eels, when pelican eels are first born, they start in the leptocephalus stage, meaning that they are extremely thin and transparent. Until they reach their juvenile stage, they interestingly have very small body organs and do not contain any red blood cells. As they mature, the males undergo a change that causes enlargement of the olfactory organs, responsible for the sense of smell, and degeneration of the teeth and jaws. The males also have defined reproductive organs. In a studied male, the testes occupied a majority of the space in the stomach cavity where the stomach had seemed to have shrunk. The females, on the other hand, remain relatively unchanged as they mature. The large olfactory organs in the sexually-mature males indicates that they may locate their mates through pheromones released by the females. Many researchers believe that the eels die shortly after reproduction. Reproducing later in life is thought to be a strategy that increases the likelihood of offspring survival for E. pelecanoides.
Distribution and habitat
The pelican eel has been found in the temperate and tropical areas of all oceans. In the North Atlantic, it seems to have a range in depth from . One Canadian-arctic specimen was found in Davis Strait at a depth of , and also across the coasts of Greenland. More recently, pelican eels have been spotted off the coast of Portugal, as well as near Hawaiian islands.
Interactions with humans
Because of the extreme depths at which it lives, most of what is known about the pelican eel comes from specimens that are inadvertently caught in deep sea fishing nets. Although once regarded as a purely deep-sea species, since 1970, hundreds of specimens have been caught by fishermen, mostly in the Atlantic Ocean. In October 2018, the first direct observation of a gulper eel was made by a group of researchers near the Azores. The team witnessed the aggressive nature of the eel's hunting process, as it was constantly moving around in the water column to attempt to find prey. In September 2018, the E/V Nautilus team also witnessed a juvenile gulper eel inflating its mouth in attempt to catch prey in the Papahānaumokuākea Marine National Monument (PMNM). Until these recent explorations, not much had been analyzed by researchers of the behavior of gulper eels.
Phylogenetic relationship to other species
In 2003, researchers from the University of Tokyo sequenced mitochondrial DNA (mtDNA) from specimens of Eurypharynx pelicanoides and Saccopharynx lavenbergi. After comparing the sequences from the specimens with other known sequences, specifically the non-coding regions, they found that E. pelicanoides and S. lavenbergi were closely related and genetically distinct from anguilliformes due to the high frequency of similarity on these regions.
| Biology and health sciences | Anguilliformes | Animals |
1888553 | https://en.wikipedia.org/wiki/Eastern%20newt | Eastern newt | Eastern newts are 2-4'inch long in length. These animals are common aquarium pets, being either collected from the wild or sold commercially. The striking bright orange juvenile stage, which is land-dwelling, is known as a red eft. Some sources blend the general name of the species and that of the red-spotted newt subspecies into the eastern red-spotted newt (although there is no "western" one).
Subspecies
The eastern newt includes these four subspecies:
Red-spotted newt (Notophthalmus viridescens viridescens)
Broken-striped newt (Notophthalmus viridescens dorsalis)
Central newt (Notophthalmus viridescens louisianensis) - Central newts measure from to in length. They are brown or green, with fine black dots all over the body. There may be a row of red spots on each side of the body. The belly is yellow or orange and is noticeably lighter than the rest of the body. The skin of newts is not as slippery as the skin of salamanders and may appear to be rough and dry for parts of their lives.
Peninsula newt (Notophthalmus viridescens piaropicola)
Life stages
Eastern newts have a lifespan of about 8–10 years in the wild, but some individuals have been known to live up to 15 years. Eastern newts have three stages of life: (1) the aquatic larva or tadpole, (2) the red eft or terrestrial juvenile stage, and (3) the aquatic adult.
Larva
The larva stage is a period of 2 to 5 months. The larva possesses gills and does not leave the pond environment where it was hatched. Larvae are brown-green, and shed their gills when they transform into the red eft. The larval Eastern Newt is the most heavily preyed upon stage. They are commonly predated on by fish, aquatic insects, and other adult newts (Brossman 2014).
Red eft
The red eft (juvenile) stage is a bright orangish-red, with darker red spots outlined in black. An eastern newt can have as many as 21 of these spots. The pattern of these spots differs among the subspecies. An eastern newt's time to get from larva to eft is about three months. During this stage, the eft may travel far, acting as a dispersal stage from one pond to another, ensuring outcrossing in the population. The striking coloration of this stage is an example of aposematism — or "warning coloration" — which is a type of antipredator adaptation in which a "warning signal" is associated with the unprofitability of a prey item (i.e., the saturation of the eft's tissues with tetrodotoxin) to potential predators. Their tetrodotoxin is a neurotoxin which is also the strongest emetic that is known. Sometimes the juvenile will continue its aquatic existence also after metamorphosis.
Adult
After two or three years, the eft finds a pond and transforms into the aquatic adult. The adult's skin is a dull olive green dorsally, with a dull yellow belly, but retains the eft's characteristic black-rimmed red spots. It develops a larger, blade-like tail and characteristically slimy skin.
It is common for the peninsula newt (N. v. piaropicola) to be neotenic, with a larva transforming directly into a sexually mature aquatic adult, never losing its external gills. The red eft stage is in these cases skipped.
Habitat
Eastern newts are at home in both coniferous and deciduous forests. Habitat preferences include shallow water, quiet stretches of streams, swamps or ditches, lakes and ponds with heavy submerged vegetation. and nearby damp woodlands. They need a moist environment with either a temporary or permanent body of water, and thrive best in a muddy environment. Eastern newts have a preference for certain types of habitats, with males preferring more open, aquatic habitats and females preferring more forested, terrestrial habitats. This preference may be related to the different roles that males and females play in the reproductive process, with males typically being more active in courtship and females spending more time on land preparing to lay eggs.
Eastern newts may travel far from their original location during the eft stage. They are most active during warm rainy periods—warmer than —and will hide under leaf litter in dry weather. Red efts may often be seen in a forest after a rainstorm. Adults prefer a muddy aquatic habitat, but will move to land during a dry spell. Eastern newts have some amount of toxins in their skin, which is brightly colored to act as a warning. Even then, only 2% of larvae make it to the eft stage. Some larvae have been found in the pitchers of the carnivorous plant Sarracenia purpurea.
Diet
Eastern newts are carnivorous, feeding on a variety of prey every two to three days. As larvae, they feed on small aquatic invertebrates, and as adults, they eat insects, worms, snails, and other small invertebrates. Eastern newts eat a variety of prey, such as insects, springtails, soil mites, small mollusks and crustaceans, young amphibians, worms, and frog eggs. They also eat a lot of snails, beetles, ants, and mosquito larvae, with an annual ingestion of about 35,000 kcal. Their dietary habits prove to be beneficial to humans because they help to control insect populations and maintain balance to their habitats. Eastern newts are a vital part of many ecosystems, serving as both predators and prey.
Behavior
Eastern newts have a number of natural predators, including fish, snakes, birds, and larger salamanders. They have several defenses against these predators, including their bright coloring, which serves as a warning signal, and their ability to secrete toxins from their skin as a defense mechanism.
Adaptability
Eastern newts are highly sensitive to changes in their environment and are able to detect and respond to changes in water quality and temperature. This sensitivity allows them to thrive in a variety of habitats, but it also makes them vulnerable to environmental changes and pollution. In fact, eastern newts are considered a sensitive species, meaning that they are often used as indicators of ecosystem health. When populations of eastern newts decline, it can be a sign of environmental stress or degradation. Newt populations are threatened by deforestation, habitat fragmentation, and pollution.
Hibernation
Eastern newts are ectothermic, relying on external sources of heat to regulate their body temperature. They are most active during the warmer months of the year, but they can also be found in more temperate climates where they may be active year-round. Eastern Newts have showcased a resistance to a wide variety of temperatures, altering their body chemistry and being able to survive and breed even under ice in winter conditions. During the winter months, some eastern newts will often burrow underground or seek shelter in logs or other debris to avoid the cold. However, studies have shown that some do not engage in hibernation, depending on the location of the species.
Homing
Eastern newts home using magnetic orientation. Their magnetoreception system seems to be a hybrid of polarity-based inclination and a sun-dependent compass. Shoreward-bound eastern newts will orient themselves quite differently under light with wavelengths around 400 nm than light with wavelengths around 600 nm, while homing newts will orient themselves the same way under both short and long wavelengths. Ferromagnetic material, probably biogenic magnetite, is likely present in the eastern newt's body.
A study determined that the home range size for Eastern newts is primarily affected by food availability, substrate humidity, but not affected by dispersal ability, competition, shelter availability, or predator avoidance. Distance traveled depended on humidity and precipitation. The mean distance traveled overnight was about 15 m, with longest trails ranging over 70 m.
Reproduction
Eastern newts breed once per year, when breeding starts in late fall until early spring. They are known to be polygynandrous, with females and males mating with multiple partners. Males have preference towards larger females, while no evidence for female preference during mating was found. The breeding migration often happens more with rainfall. The male's spots attract females, luring them to him with fanning motions of his tail, causing a pheromone to be released. Once the female has chosen a mate, the male will deposit a spermatophore, a package of sperm, onto the ground, which the female will then pick up and fertilize her eggs with. The female will lay her eggs in the water, attaching them to submerged vegetation or other objects. 200–400 eggs are laid in a single batch, with incubation period of 3–8 weeks. For the normal and healthy development of gonads, fat-bodies are needed in proximity of the developing organs to ensure proper reproduction ability.
Social interactions
The behavior of eastern newts is also influenced by their social interactions with other members of their species. Eastern newts exhibit social hierarchy, with dominant individuals exhibiting aggressive behaviors towards subordinates. This social hierarchy is thought to be related to the distribution of resources, with dominant individuals having access to more food and better mating opportunities. One such behavior is territoriality, where individuals will defend a specific area or resource from other members of their species. This behavior is commonly seen in males during the breeding season, when they will defend a territory in order to attract females and ensure access to mating opportunities.
Survival advantages
Secretion of toxins through the skin protects the newt from predators, and should therefore not be handled with bare hands. The red colors of the adult newt also act as a warning sign for predators. Its ventral surface has poison glands, which makes predators reluctant to eat it. However, one study observed a Belted Kingfisher (Megaceryle alcyon) beat an eastern newt on a nest box 15 times before eating it. This special toxin is known as tetrodotoxin. Several studies have found that newt larvae increase the production of this toxin while in the presence of predators (dragonflies). Tetrodotoxin is known to cause muscle paralysis, skin irritation, and even death in predators, although some mantis species have shown a resilience to this toxin, and predatory sunfish are not deterred by the toxin. The Eastern newt also has a greater tail depth and is capable of swimming quickly away from aquatic predators.
Limb regeneration
Eastern newts are able to regenerate their limbs that were lost to an injury. Forelimb regeneration has been considered to be close to the forelimb development; genes that play a role in forelimb regeneration are known to also be expressed in its developmental stages. In addition, they are capable of regenerating their spinal cord, heart, and other organs. This ability is thought to be related to their high levels of stem cells, which allow them to repair and regenerate damaged tissues.
Conservation concerns
Although eastern newts are widespread throughout North America, they, like many other species of amphibians, are increasingly threatened by several factors including habitat fragmentation, climate change, invasive species, over-exploitation, and emergent infectious diseases. The biodiversity of amphibians across the United States is considered to be threatened due to the loss of wetlands and furthermore, their connectivity; since the 1780s, more than 53% of wetlands in the United States have been lost. For example, a study found the toxicity of coal-tar pavement on eastern newts sublethal, decreasing their righting ability and swimming speed. Wild eastern newts are known hosts of Batrachochytrium dendrobatidis and Ranavirus, as well as the mesomycetozoan Amphibiocystidium ranae. They are also highly susceptible to the newly emergent chytrid fungus Batrachochytrium salamandrivorans.
Gallery
| Biology and health sciences | Salamanders and newts | Animals |
1888817 | https://en.wikipedia.org/wiki/Central%20force | Central force | In classical mechanics, a central force on an object is a force that is directed towards or away from a point called center of force.
where is the force, F is a vector valued force function, F is a scalar valued force function, r is the position vector, ||r|| is its length, and is the corresponding unit vector.
Not all central force fields are conservative or spherically symmetric. However, a central force is conservative if and only if it is spherically symmetric or rotationally invariant.
Properties
Central forces that are conservative can always be expressed as the negative gradient of a potential energy:
(the upper bound of integration is arbitrary, as the potential is defined up to an additive constant).
In a conservative field, the total mechanical energy (kinetic and potential) is conserved:
(where 'ṙ' denotes the derivative of 'r' with respect to time, that is the velocity,'I' denotes moment of inertia of that body and 'ω' denotes angular velocity), and in a central force field, so is the angular momentum:
because the torque exerted by the force is zero. As a consequence, the body moves on the plane perpendicular to the angular momentum vector and containing the origin, and obeys Kepler's second law. (If the angular momentum is zero, the body moves along the line joining it with the origin.)
It can also be shown that an object that moves under the influence of any central force obeys Kepler's second law. However, the first and third laws depend on the inverse-square nature of Newton's law of universal gravitation and do not hold in general for other central forces.
As a consequence of being conservative, these specific central force fields are irrotational, that is, its curl is zero, except at the origin:
Examples
Gravitational force and Coulomb force are two familiar examples with being proportional to 1/r2 only. An object in such a force field with negative (corresponding to an attractive force) obeys Kepler's laws of planetary motion.
The force field of a spatial harmonic oscillator is central with proportional to r only and negative.
By Bertrand's theorem, these two, and , are the only possible central force fields where all bounded orbits are stable closed orbits. However, there exist other force fields, which have some closed orbits.
| Physical sciences | Classical mechanics | Physics |
267008 | https://en.wikipedia.org/wiki/MTR | MTR | The Mass Transit Railway (MTR) is a major public transport network serving Hong Kong. Operated by the MTR Corporation (MTRCL), it consists of heavy rail, light rail, and feeder bus services, centred around a 10-line rapid transit network, serving the urbanised areas of Hong Kong Island, Kowloon, and the New Territories. The system encompasses of railways, as of December 2022, with 179 stations—including 99 heavy rail stations, 68 light rail stops and 1 high-speed rail terminus.
Under the government's rail-led transport policy, the MTR system is a common mode of public transport in Hong Kong, with over five and a half million trips made on an average weekday consistently achieving a 99.9% punctuality rate on its arrivals and departures. As of 2018, the MTR has a 49.3% share of the franchised public transport market, making it the most popular transport option in Hong Kong. The integration of the Octopus smart card fare-payment technology into the MTR system in September 1997 has further enhanced the ease of commuting.
History
Initial proposals
During the 1960s, the government of Hong Kong saw a need to accommodate increasing road traffic as Hong Kong's economy grew rapidly. In 1966, British transport consultants Freeman, Fox, Wilbur Smith & Associates were appointed to study the transport system of Hong Kong. The study was based on the projection of the population of Hong Kong for 1986, estimated at 6,868,000. On 1 September 1967, the consultants submitted the Hong Kong Mass Transport Study to the government, which recommended the construction of a rapid transit rail system in Hong Kong. The study suggested that four rail lines be developed in six stages, with a completion date set between December 1973 and December 1984. Detailed locations of lines and stations were presented in the study. These four lines were the Kwun Tong line (from Mong Kok to Ma Yau Tong), Tsuen Wan line (from Admiralty to Tsuen Wan), Island line (from Kennedy to Chai Wan Central), and Shatin line (from Tsim Sha Tsui to Wo Liu Hang).
The study was submitted to the Legislative Council on 14 February 1968. The consultants received new data from the 1966 by-census on 6 March 1968. A short supplementary report was submitted on 22 March 1968 and amended in June 1968. The by-census indicated that the projected 1986 population was reduced by more than one million from the previous estimate to 5,647,000. The dramatic reduction affected town planning. The population distribution was largely different from the original study. The projected 1986 populations of Castle Peak New Town, Sha Tin New Town, and, to a lesser extent, Tsuen Wan New Town, were revised downwards, and the plan for a new town in Tseung Kwan O was shelved. In this updated scenario, the consultants reduced the scale of the recommended system. The supplementary report stated that the originally suggested four tracks between Admiralty station and Mong Kok station should be reduced to two, and only parts of the Island line, Tsuen Wan line, and Kwun Tong line should be constructed for the initial system. The other lines would be placed in the list of extensions. This report led to the final study in 1970.
In 1970, a revised system with four lines was laid out in the British consultants' new report, Hong Kong Mass Transit: Further Studies. The four lines were to be the Kwun Tong line, Tsuen Wan line, Island line, and East Kowloon line. The lines that were eventually constructed were somewhat different compared to those presented in this report and the Hong Kong Mass Transport Study.
In 1972, the Hong Kong government authorised construction of the Initial System, a system that roughly translates to today's Kwun Tong line between Kwun Tong and Prince Edward, Tsuen Wan line between Mei Foo and Admiralty, and Island line between Sheung Wan and Admiralty. The Mass Transit Steering Committee, chaired by the Financial Secretary Philip Haddon-Cave, began negotiations with four major construction consortia in 1973. The government's intention was to tender the entire project, based on the British design, as a single tender at a fixed price. A consortium from Japan, led by Mitsubishi, submitted the only proposal within the government's $5-billion price ceiling. They signed an agreement to construct the system in early 1974, but in December of the same year, pulled out of the agreement for reasons stemming from fears of the oil crisis.
Modified initial system
Several weeks later, in early 1975, the Mass Transit Steering Group was replaced by the Mass Transport Provisional Authority, which held more executive powers. It announced that the Initial System would be reduced to and renamed the "Modified Initial System" (now part of the Kwun Tong and Tsuen Wan lines). Plans for a single contract were abandoned in favour of 25 engineering contracts and 10 electrical and mechanical contracts. On 7 May 1975 the Legislative Council passed legislation setting up the government-owned Mass Transit Railway Corporation (MTRC) to replace the Mass Transport Provisional Authority, the Mass Transit Railway Ordinance.
Construction of the system began on 11 November 1975. The northern section was completed on 30 September 1979 and was opened on 1 October 1979 by Governor Murray MacLehose. Trains on this route ran from Shek Kip Mei to Kwun Tong in Phase 1, Tsim Sha Tsui to Kwun Tong in Phase 2 in December 1979, and Chater to Kwun Tong in the last phase, initially in a four-car configuration. The first train drivers were trained on the London Underground. It was designed by a consortium of consultants led by Freeman Fox and Partners. On later extensions to the railway the stations were designed under the supervision of Roland Paoletti, the chief architect at MTR.
The full Modified Initial System was opened on 12 February 1980 by Princess Alexandra, who rode the inaugural train through the immersed tube beneath Victoria Harbour to Central station. Trains were gradually extended to six cars to accommodate an increase in passenger numbers.
Line extensions
The government approved construction of the Tsuen Wan line in 1977, then known as the Tsuen Wan Extension, and works commenced in November 1978. The project added a section to the MTR system, from Prince Edward station to Tsuen Wan. The line started service on 17 May 1982 with a total cost of construction (not adjusted for inflation) of HK$4.1 billion (US$526 million). The plan was modified from that in the 1970 report Hong Kong Mass Transit: Further Studies, with Kwai Chung station, Lap Sap Wan station, and a planned depot in Kwai Chung next to Lap Sap Wan station being replaced by stations in Kwai Hing and Kwai Fong and a depot in Tsuen Wan. Several stations also had names different to that during planning: So Uk station became Cheung Sha Wan, Cheung Sha Wan became Lai Chi Kok, and Lai Chi Kok became Lai Wan (later renamed Mei Foo).
When service of this line started, the section of the Kwun Tong line from Chater to Argyle (since renamed Central and Mong Kok respectively) was transferred to the Tsuen Wan line. Thus, Waterloo station (since renamed Yau Ma Tei) became the terminus of the Kwun Tong line, and both Argyle and Prince Edward stations became interchange stations. This change was made because system planners expected the patronage of the Tsuen Wan line to exceed that of the Kwun Tong line. This forecast proved to be accurate, necessitating a bypass from the northwestern New Territories to Hong Kong Island. The Tung Chung line was therefore launched in 1998 with an interchange station at Lai King for that purpose.
Although land acquisitions were made for a station at Tsuen Wan West (near Tsuen King Circuit), beyond Tsuen Wan station, as part of the Tsuen Wan branch, the station was never built. This is not to be confused with the modern-day Tsuen Wan West station on Tuen Ma line, which lies on a newly reclaimed area near the former ferry pier.
Since opening in 1982, the Tsuen Wan line is the line whose alignment has remained the same for the longest time. For example, the Kwun Tong line's alignment has changed three times since its opening—the taking over of Tsuen Wan line from Mong Kok to Central, the taking over of Eastern Harbour Crossing section by the Tseung Kwan O line, and its extension to Whampoa.
Government approvals were granted for construction of the Island line in December 1980. Construction commenced in October 1981. On 31 May 1985, the Island line was opened with service between Admiralty station and Chai Wan station. Both Admiralty and Central stations became interchange stations with the Tsuen Wan line. Furthermore, each train was extended to eight cars. On 23 May 1986, the Island line was extended to Sheung Wan station. Construction was delayed for one year, as government offices which were located over the station had to be moved before the construction could start.
In 1984, the government approved the construction of the Eastern Harbour Crossing, a tunnel to be used by cars and MTR trains. The Kwun Tong line was extended across the harbour on 5 August 1989 to Quarry Bay station, which became an interchange station for the Kwun Tong line and the Island line. An intermediate station, Lam Tin, started operations on 1 October 1989.
Airport connection
The decision was made in October 1989 to construct a new international airport at Chek Lap Kok on Lantau Island to replace the overcrowded Kai Tak International Airport. The government invited the MTRC to build a train line, then known as the Lantau Airport Railway, to the airport. Construction started in November 1994, after the Chinese and British governments settled their financial and land disagreements.
The new line was included in the financing plans of the new Hong Kong International Airport as the airport was not considered viable without direct public transport links. Construction costs were also shared by the MTRC, which was granted many large-scale developments in the construction plans for the new stations.
The Lantau Airport Railway included two MTR lines, the Tung Chung line and the Airport Express. The Tung Chung line was officially opened on 21 June 1998 by Hong Kong Chief Executive Tung Chee-hwa, and service commenced the next day. The Airport Express opened for service on 6 July 1998 along with the new Hong Kong International Airport.
The Airport Express also offers flight check-in facilities at Kowloon station and Hong Kong station—the in-town check-ins offer a more convenient and time-saving routine; a free shuttle bus service transports travellers from these stations to their respective hotels as well. Porters are also available to help transport luggage from and onto trains. It is the second most popular means of transport to the airport after buses. In 2012, it had a 21.8 per cent of share of the traffic to and from the airport. However, this has declined from a peak of 32 per cent in 1999.
Tseung Kwan O line
The Quarry Bay Congestion Relief Works extended the Hong Kong Island end of the Kwun Tong line from Quarry Bay to North Point via a pair of tunnels. The project was initiated due to overcrowding at Quarry Bay and persistent passenger complaints about the five-minute walk from the Island line platforms to the Kwun Tong line platform. Construction began in September 1997 and was completed in September 2001 at a cost of HK$3.1 billion. As with most earlier interchange stations, a cross-platform interchange arrangement was provided here in both directions.
Construction of the Tseung Kwan O line (called the Tseung Kwan O extension line in the planning stage) was approved on 18 August 1998 to serve the growing Tseung Kwan O New Town. Construction began on 24 April 1999 and the line officially opened in 2002. It took over the existing Kwun Tong line tracks running through the Eastern Harbour Tunnel, so that the full line stretches from Po Lam to North Point. When the line opened, the Kwun Tong line was extended to Tiu Keng Leng on the new line. Construction costs were partly covered by the Hong Kong Government and private developers which linked construction of the Tseung Kwan O line to new real estate and commercial developments.
Interchange stations
The interchange between the Kwun Tong line and the Tsuen Wan line (except Yau Ma Tei) as well as that between the Kwun Tong line and the Tseung Kwan O line, are two stations long, allowing cross-platform interchange wherein a passenger leaves a train on one side of the platform and boards trains on the other side of the platform for another line. For example, when passengers are travelling on the Kwun Tong line towards Tiu Keng Leng, getting off at Yau Tong would allow them to switch trains across the platform for the Tseung Kwan O line towards North Point. Whereas, staying on the train and reaching Tiu Keng Leng would allow them to board the Tseung Kwan O line trains towards Po Lam/LOHAS Park. This design makes interchanging more convenient and passengers do not have the need to change to different levels. However this interchange arrangement is not available for all transferring passengers at Kowloon Tong, Central, Hong Kong, Quarry Bay, Nam Cheong (except transfer between Tuen Mun and Hong Kong bound trains), Mei Foo, Tai Wai (only between southbound Tuen Ma line and East Rail line trains) and Sunny Bay (except transfer between Tung Chung and Disneyland Resort bound trains) stations, mainly because this service is available only when there are two continuous stations shared as interchange stations by two lines.
Two major works were undertaken to ease interchange between the Kwun Tong line and East Rail line. The modification of Kowloon Tong station started in June 2001. A new pedestrian link to Kowloon Tong station southern concourse and a new entrance (Exit D) opened on 15 April 2004 to cope with the increase in interchange passenger flow. Modification to Tsim Sha Tsui station involved upgrading station facilities and concourse layout to facilitate access from the East Tsim Sha Tsui station via its pedestrian links. New entrances to the subway links were opened on 19 September 2004 (Exit G) and 30 March 2005 (Exit F), with the whole scheme completed in May 2005.
Disneyland Resort line
The Disneyland Resort line, previously known as Penny's Bay Rail Link, provides service to the Hong Kong Disneyland Resort station which was opened on 12 September 2005. Services to Sunny Bay station on the Tung Chung line started in 1 June 2005, but it was only opened to staff of Disneyland at first. It was finally opened to the general public two months later, on 8 August 2005. The new line and the Disneyland Resort station opened on 1 August 2005. It is a single-track railway that runs between Sunny Bay station and Disneyland Resort station. The Disneyland Resort station itself was designed to blend in with the ambiance of the resort. The line operates fully automated trains running every four to ten minutes without a driver. The carriages are refurbished M-train rolling stock to match the recreational and adventurous nature of the 3.5-minute journey.
Airport Express extension
The AsiaWorld–Expo station is an extension of the Airport Express serving a new international exhibition centre, AsiaWorld–Expo, at Hong Kong International Airport. The station opened on 20 December 2005 along with the exhibition centre. To cope with the projected increase in patronage, Airport Express trains were lengthened to eight cars from the previous seven. Additional trains are also deployed on the Tung Chung line during major exhibitions and events.
Partial privatisation and merger
On 5 October 2000 the operator of the MTR network, the Mass Transit Railway Corporation (MTRC), became Hong Kong's first rail company to be partially privatised, marking the beginning of the Hong Kong government's initiative to reduce its interests in public utilities. Prior to its listing on the Hong Kong Stock Exchange, the Mass Transit Railway Corporation (MTRC) was wholly owned by the Hong Kong government. The offering involved the sale of about one billion shares, and the company now has the largest shareholder base of any company listed in Hong Kong. In June 2001, MTRCL was transferred to the Hang Seng Index.
MTRCL has often developed properties next to stations to complement its profitable railway business. Many recently built stations were incorporated into large housing estates or shopping complexes. For example, Tsing Yi station is built next to the Maritime Square shopping centre and directly underneath the Tierra Verde housing estate.
On 11 April 2006, MTRCL signed a non-binding memorandum of understanding with the Hong Kong government, the owner of Kowloon–Canton Railway Corporation, to merge the operation of the two railway networks in Hong Kong in spite of the strong opposition of KCRC staff. The minority shareholders of the corporation approved the proposal at an extraordinary general meeting on 9 October 2007, allowing MTRCL to take over the operation of the KCR network and combine the fare system of the two networks on 2 December 2007.
On 2 December 2007 the Kowloon-Canton Railway Corporation (KCRC) granted a 50-year service concession (which may be extended) of the KCR network to MTRCL, in return for making annual payments to KCRC, thereby merging the railway operations of the two corporations under MTRCL's management. At the same time MTRCL changed its Chinese name from "地下鐵路有限公司" (Subway Limited Company) to "香港鐵路有限公司" (Hong Kong Railway Limited Company), but left its English name unchanged; at the same time the system's Chinese name changed from "地鐵" ("underground railway") to "港鐵" ("Hong Kong Railway"). After the merger, the MTR network included three more lines—East Rail line, West Rail line, and Ma On Shan line (now the Tuen Ma line)—as well as the light rail network and Guangdong through train to Guangzhou.
On 28 September 2008, fare zones of all urban lines, East Rail line, Ma On Shan line, and West Rail line were merged. A passenger could travel on these networks with only one ticket, except where a transfer is made between Tsim Sha Tsui and East Tsim Sha Tsui stations, where two tickets are required. Student discounts on Octopus Card were also issued.
Recent extensions
The MTR system has been extended numerous times since the railway merger. Relevant projects include the LOHAS Park spur line (2009), the Kowloon Southern Link (2009), the West Island line (2014), the Kwun Tong line extension (2016), the South Island line (2016), Tuen Ma line Phase 1 (2020) and Phase 2 (2021) and the East Rail line extension (2022).
The LOHAS Park Spur Line is an extension of the Tseung Kwan O line, splitting off after Tseung Kwan O station. It serves the new residential development of LOHAS Park (formerly "Dream City"), a estate with fifty residential towers. The project is divided into 9 to 13 phases and is about halfway complete as of 2016. These high rises sit above LOHAS Park station, which opened on 26 July 2009.
The West Island line, first put forward to the government on 21 January 2003, is an extension of the Island line. It serves the Western District of Hong Kong Island. The construction of the West Island line started on 10 August 2009. Kennedy Town station and HKU station opened on 28 December 2014. Sai Ying Pun station opened later, on 29 March 2015, due to construction delays.
A proposal to extend the existing Kwun Tong line to Whampoa Garden was made in April 2006 and approved in March 2008 as part of the bid for the Sha Tin to Central Link. Two new stations at Whampoa and Ho Man Tin opened on 23 October 2016.
The South Island line opened on 28 December 2016 between Admiralty and South Horizons, linking the MTR to Southern District for the first time. With the opening of the South Island line, all 18 districts of Hong Kong are served by the MTR.
The first section of the Tuen Ma line, an extension of the former Ma On Shan line connecting Tai Wai via Hin Keng and Diamond Hill to Kai Tak station, opened on 14 February 2020. The second and final section of the line was completed and opened on 27 June 2021, linking the previously opened Tuen Ma Line Phase One and the West Rail Line together connecting from Kai Tak station to Hung Hom station.
An extension of the East Rail line, phase two of the Sha Tin to Central Link (SCL) from Hung Hom station to Admiralty station across Victoria Harbour was completed and opened on 15 May 2022. An intermediate station was opened at Exhibition Centre.
Future extensions
Tuen Mun South extension
The Tuen Mun South extension on the Tuen Ma line is a proposed extension to a new western terminus, Tuen Mun South, near Tuen Mun Ferry Pier. The extension will extend the line southwards from the current terminus at Tuen Mun station. It will include the construction of the A16 station (placeholder name used by MTR) and the new terminus Tuen Mun South station.
Additionally, the addition of a new infill station, Hung Shui Kiu station, along the Tuen Ma line between Siu Hong station and Tin Shui Wai station is currently under planning. It may be built depending on the development of the Hung Shui Kiu New Town.
Northern Link
The Northern Link is a proposed new line which connects Tuen Ma line with the Lok Ma Chau Spur Line of East Rail line. It also has Au Tau, Ngau Tam Mei, San Tin, a future interchange station between East Rail line and Northern Link, Kwu Tung, which will become a terminus for Northern link. This line would serve the future Northern Metropolis (which is in current planning) by the Hong Kong government. It would help to connect planned population centres isolated in the New Territories with Kowloon and Hong Kong.
Construction of Kwu Tung station began in 29 September 2023, and is expected to be completed in 2027, while construction of the Northern Link is expected to begin in 2025 and is scheduled to commence service in 2034.
Tung Chung line extension
The Tung Chung line extension will extend the Tung Chung line to the west by approximately 1.3 kilometres. Two new stations will also be built, namely; Tung Chung West and Tung Chung East, with Tung Chung West serving as the new terminus of the Tung Chung line. Construction began on 25 May 2023 and is expected to be completed in 2029. The Oyster Bay station is a planned infill station between Sunny Bay station and the future Tung Chung East station. It is expected to be complete in 2030.
East Kowloon line
The East Kowloon line is planned to serve the East Kowloon area to Tseung Kwan O New Town via the hilly Sau Mau Ping residential area.
South Island line (West)
The South Island line (West) was part of the same original proposal as the South Island line, and would connect HKU to Wong Chuk Hang around the west coast of Hong Kong Island, however, construction has not started .
North Island line
The North Island line is a planned extension of the Tseung Kwan O line that will interchange at the future Tamar station with the Tung Chung line. It will alleviate traffic in the northern part of Hong Kong Island. There will be three new stations: Tamar, Exhibition Centre (which will be an interchange between the North Island line and the North South Corridor), and Causeway Bay North. There is currently no proposed construction time for this line, however in the original proposal, construction was expected to begin in 2026 and commence service by 2040. The cost is estimated to be HK$20 billion in 2013 prices.
Infrastructure
Rail network
Rolling stock
Eight types of electric multiple unit rolling stock operate on the MTR network and five generations of light rail vehicles operate on the light rail network. All use either rail gauge (near standard gauge) or (standard gauge). Except for Airport Express trains, all trains are designed to cope with high patronage, for example, through seating arrangements, wider train cars, additional ventilation fans, brighter train lighting, and additional sets of extra-wide doors. These configurations allow the MTR to run at up to 101,000 passengers per hour per direction (p/h/d) on its busy suburban East Rail line and around 85,000 p/h/d on its urban metro network.
Former rolling stock
Metro Cammell EMU (DC)
Known as M-trains, these are the oldest model of train in operation. M-Trains can be divided into different "stocks". The M-stock (or "CM-stock") of M-Train are the oldest trains on the MTR, built originally by Metro-Cammell (now Alstom) and refurbished by United Goninan. The M-train uses sliding doors, unlike K-stocks and Grupo CAF Trains which use plug doors. They are in service on Kwun Tong line, Tsuen Wan line, Island line and Tseung Kwan O line.
The Disneyland Resort line uses driverless M-trains with their appearance overhauled to suit the atmosphere and theme of the line. Windows on each carriage and the handrails inside are made into the shape of Mickey Mouse's head, and there are bronze-made Disney characters decorating the interior of the carriages.
All 93 sets of the M-trains will be retired from service by 2030, and be replaced by the Q-Trains.
Adtranz-CAF EMU
The Tung Chung line and the Airport Express use CAF Trains tailored to their respective lines. Initially run in seven-car formations, they have now been lengthened to eight cars. These two variations are built jointly by Adtranz (now Bombardier Transportation) and Grupo CAF (CAF) between 1994 and 1997. Since 2006, K-stock has also been used on the Tung Chung line.
Rotem EMU (K-trains)
The K-stock was built jointly by Mitsubishi Heavy Industries and Hyundai Rotem and first put into service on the Kwun Tong line. These were eventually transferred to the Tseung Kwan O line in 2009. Subsequently, in 2006, four additional sets joined the Tung Chung line to cope with the increasing passenger traffic. K-Stock trains have come under criticism when they were first put into service due to delays and door safety issues. Along with other service reliability issues, there have been incidents where passengers have been injured by its doors, leading to the MTRCL "minimising the number of Korean trains for passenger service until a higher reliability of the systems concerned is achieved".
CNR Changchun EMU (C-trains)
A contract (C6554-07E) for 10 new sets of trains was awarded to Changchun Railway Vehicles Co. Limited in October 2008 with a further 12 trains ordered in the summer of 2011. These were delivered to Hong Kong between 2011 and 2013 to enhance train frequency on the existing lines to cater for increased patronage on the Island, Kwun Tong, Tsuen Wan and Tseung Kwan O lines.
These trains feature new 22" LCD TVs, like their counterparts on former KCR lines, and as a result are equipped with MTR In-Train TV, offering infotainment such as news and announcements. The first of these trains entered revenue service on 7 December 2011 on the Kwun Tong line.
The South Island line uses a similar train known as the S-train. Unlike the original C-train, the pantographs are on the "A" cars. All cars are powered, so there are no trailer cars. The S-train is also only three cars long and are driverless, with the driver's cab removed to make standing space for passengers. However, the S-trains can still be operated manually in the event of an emergency.
SP1900/1950 EMU (IKK-trains)
The Tuen Ma line uses the SP1900, also known as the IKK (Itochu-Kinki-Kawasaki) train. The electrification system used on Tuen Ma line is , as opposed to the used on the urban lines. Should the need arise in the future, dual voltage trains such as those used on Oresund Bridge would be required.
The rolling stock is from the former KCRC network (KCR East Rail, West Rail and Ma On Shan Rail). They did not receive major changes after the merger of the two companies except for the updated route map, the exterior company logo and such. The capability of this EMU fleet is similar to those on the urban network. Starting from 2015, the West Rail and Ma On Shan line trains have been lengthened to 8 cars while the East Rail line 12 car sets have been withdrawn from the line in 2021. All train sets will receive larger TV displays and dynamic route map displays above every door, and will run on the Tuen Ma line in the future.
Hyundai Rotem EMU (R-stock)
In December 2012, the MTRC announced that new contracts had been awarded to Hyundai Rotem for 37 new nine-car trains to be used on the Sha Tin to Central Link. These trains are expected to replace the Metro Cammell EMUs that currently run on the East Rail line. The new R-stock trains are wider than existing units and can accommodate more passengers per car; however, the length of each train will be shortened from the current 12-car configuration used on the Metro Cammell and SP1900 EMUs to nine cars. This is due to space constraints imposed by new underground platforms on the Sha Tin to Central Link. MTR will also upgrade existing signalling systems used on the East Rail line which will enable trains to operate at two-minute headways on average, instead of the current three-minute interval, which the MTRC expects will be able to compensate for the loss of capacity resulting from the shorter trains. However, there are concerns from local residents that this will not be effective.
Light rail vehicles
Light rail rolling stock were ordered from four different manufacturers: Commonwealth Engineering (Comeng), Kawasaki Heavy Industries, United Goninan and CRRC Nanjing Puzhen. They are designed to run on the standard gauge and use delivered through overhead lines. Trains comprise one or two carriages, where the second carriage functions as only a trailer. The arrangement allows each car to carry approximately 300 passengers with 26 seats, while four sets of poach seats provide flexible riding for passengers.
The light rail trains are being modernised as part of a 20th anniversary activity. Trains will include better disabled facilities as well as a totally new interior. The MTR will refurbish 69 older trains (revised to 68 as one was scrapped following a traffic incident) and buy 22 new ones. The first trains have been completed and were scheduled to be put into service in November 2009. The whole project is expected to be completed in 2011.
Another batch of 40 Phase V trains have been ordered from CRRC Nanjing Puzhen, which will replace 30 Kawasaki Phase II trains which will not be refurbished, and the 10 additional trains will be used for enhancing services. The first pair of trains (leading + trailer car) entered service in 2020.
Maintenance trains
In addition to the passenger electric multiple units as covered above, MTR also uses various types of work trains for maintenance purposes:
On the urban lines and Lantau Airport Railway, several different battery-electric locomotives and diesel locomotives are used to haul various work trains (including ultrasonic test vehicles and specialist wagons used for overhead wiring access, cable laying, rail transport, tunnel repair etc.), with the former being built by Brush Traction and the latter by Schöma. Like the passenger trains, they are also equipped with BSI couplers, albeit without the electrical connections. The locomotives were delivered in different phases:
The battery electrics were delivered in three different phases. The first phase comprises five units numbered L51 to L55 which entered service in 1983 and are capable of hauling a maximum trailing load of 109 tonnes on a 3% of gradient. The second phase has six units numbered L56 to L61 which entered service in 1989. They have the same performance as the phase 1 units and can be assumed to be of the same general design. The pantograph is located atop the cab roof, with no sign of an air conditioner. The third phase comprises twenty units numbered L62 to L82 which entered service in 1996; most are used on the Lantau Airport Railway, but some on the urban lines. They are more powerful than the earlier units, being able to haul 160 tonnes. Their design is also slightly different to phase 2 units, with the pantograph at the opposite end, and an air conditioner above the cab.
The Schöma diesel fleet came before the battery electrics. Capable of hauling 100 tonnes, they are presumed to have assisted in the construction of the network and came in three different phases: phase 1 (nine units numbered L11 to L19) in 1977, phase 2 (nine units numbered L20 to 28) in 1979 and phase 3 (eight units numbered L31 to L38) in 1983. They are presumed to be capable of multiple unit operation.
On the ex-KCR lines, diesel locomotives such as the Eurorunner that were formerly used for freight services are now used on works trains and transferring EMUs trains between depots. Maintenance trains for the East Rail line are also largely stabled at Fo Tan depot.
In 2015, MTR procured two new SF02T-FS rail milling machines from Linsinger. Instead of rail grinding, these machines mill the rail head, providing a more accurate profile and a higher quality processed surface. However, the machines are not without flaws and one of them caught fire on the Tsuen Kwan O Line on 9 May 2016 during maintenance hours.
Station facilities, amenities and services
The architecture of MTR stations is less artistic, instead focusing on structural practicability. With the high level of daily passenger traffic, facilities of the MTR stations are built with durability and accessibility in mind. After extensive retrofitting, the MTR system has become, in general, disabled-friendly—the trains have dedicated wheelchair space, the stations have special floor tiles to guide the blind safely on the platforms, and there are extra wide entry and exit gates for wheelchairs as well. Portable ramp for wheelchair users are available for boarding and alighting trains. On board the rolling stock, there are also flashing system maps on select trains while Active Line Diagrams and traditional route maps are installed on the others. Infopanels as well as on MTR In-Train TV onboard trains display important messages such as next station announcements as well as operational messages.
Telecommunications
LTE (4G) and 5G mobile phone network is in place throughout the whole MTR system of stations and tunnels allowing passengers to stay connected to the internet underground. Currently, there is full 5G network coverage in all stations and tunnels (except certain underground sections on the East Rail line and Tuen Ma line) for the MTR system has been provided by 3 Hong Kong, SmarTone and PCCW. Passengers are able to use high-speed internet on their mobile phones regardless whether the train is above ground or under ground. The MTR has already extended the Wi-Fi service to all of the Airport Express trains and the expansion of the service to other MTR routes is still under consideration by MTR. All 99 stations on the MTR offer free Wi-Fi service with a limit of 15 minutes per session and a maximum of five sessions per day.
In late 2015 it was announced that all 400 payphones in the MTR system would be removed in early 2016. The contract with the service provider, Shinetown Telecom, was expiring, and the MTR Corporation said that no one had tendered a proposal to take over the contract. As a result, the MTR system no longer has payphones.
Announcements
When the system opened, public announcements were made in British English and Cantonese by train captains and station staff. In 1992, the announcements were standardised, pre-recorded by RTHK presenter Cheri Chan Yu-yan (), who is now an assistant professor of English-language education at the University of Hong Kong, and who remains the voice of the MTR today. Since 2004, to accommodate Mainland Chinese visitors under the Individual Visit Scheme, Mandarin Chinese was added to the repertory.
Public toilets
Unlike many other metro systems around the world, "main line" MTR stations originally did not have toilet facilities available for public use. Passengers may use MTR staff toilets at all stations on request. In 2006, MTRCL said it would not consider retrofitting existing underground toilets, because of the challenge of installing new piping and toilet facilities. Only stations on the Airport Express, Tung Chung line and Disneyland Resort line had access to toilet facilities. All former KCR stations (on the East Rail line and Tuen Ma line), merged into the MTR network in 2007, have public toilets.
During the Legislative Council rail merger bill discussions, the MTR Corporation was criticised by legislators for their unwillingness to install toilets in main line stations. MTRCL indicated that it would carry out a review of the feasibility of installing public toilets at or in the vicinity of its above-ground railway stations. Discussions between the Government and MTRCL have taken into account LegCo members' request for a stronger commitment by the corporation to the provision of public toilets on new railway lines. This resulted in MTRCL agreeing to include the provision of toilets within, or adjacent to, stations in the overall design parameters for all future new railway lines, subject to planning and regulatory approval and any concerns raised by residents in the vicinity about the location of external ventilation exhausts.
Toilets have since been retrofitted into several existing MTR stations, including Sheung Wan station, Ngau Tau Kok station, Quarry Bay station, Mong Kok station, Prince Edward station, and Admiralty station. In addition, newly opened stations such as those of the West Island line have toilets. The MTR has installed public toilets at all interchanges as of 2023.
In late 2017, the MTR introduced breastfeeding rooms at 20 interchange stations. The rooms are located in back of house areas, and are available upon request to MTR staff.
Commerce and journals
Prior to the privatisation of MTRC, MTR stations only had branches of the Hang Seng Bank, Maxim's Cakes stores, and a handful of other shops. Since then, the number and types of shops have increased at stations has increased, turning some of them into miniature shopping centres. ATMs and convenience stores are now commonplace.
The MTR has contracted with publishers for the distribution of free magazines and newspapers in MTR stations. Recruit was the first free magazine which was solely distributed in stations (before railway merger) since July 1992, but the contract was terminated in July 2002. Another recruitment magazine Jiu Jik (招職), published by South China Morning Post, replaced Recruit as the only free recruitment magazine distributed in MTR stations bi-weekly. The Metropolis Daily (都市日報), published by Metro International, is the first free newspaper distributed free in MTR stations during weekdays (except public holidays); and in 2005, there is another weekend newspaper Express Post (快線週報), distributed every Saturday except public holidays. The Metropop (都市流行), a weekly magazine featuring cultural affairs and city trends also published by Metro International, started its distribution in MTR stations every Thursday since 27 April 2006, a few months after the termination of Hui Kai Guide (去街 Guide) in 2006. MTR Stations on ex-KCR lines feature two free Chinese-language newspapers, namely am730 and Headline Daily. MTR promotes reading of these newspapers by adding special coupons and promotion offers inside the newspapers, for example, a free trip to Lok Ma Chau or a free keyring. On the Kwun Tong line, East Rail line and Tuen Ma line, MTR In-Train TV is available.
MTR Bus
At various stations of the MTR network, the MTRCL (which took over from KCR) operates feeder buses which enhance the convenience of taking the MTR. These bus routes, which normally consist of one to two stops, terminate at housing estates and go past major landmarks. The feeder bus routes on the East Rail line are run under the MTR name but are operated by Kowloon Motor Bus.
Signalling
Throughout its history, MTR has used different signalling systems for its lines. The main Operations Control Centre for the entire network is located at Tsing Yi. Previous control centres were located at Fo Tan and Kam Tin for the East Rail line, Ma On Shan line and West Rail line.
On the pre-merger MTR network, wayside signals are simple two-aspect signals whose colours are namely red for "stop" and blue for "proceed according to ATO"; this is made possible by the use of automatic train operation (ATO) which provides the onboard equipment the permitted speeds via undercarriage antennas located underneath the cab whereas the signals having been sent by radio transmitters located between the rails. An automatic train protection (ATP) is also used to enforce safety.
In 1998, transmission balise-locomotive (TBL) was implemented on the East Rail line to monitor train safety. Subsequently, in 2002, ATO was also implemented on the East Rail line. However, the original British-style Automatic Warning System is still retained for use by Intercity-Through Trains. On the other hand, the Tuen Ma line uses a SelTrac moving block communications-based train control (CBTC) system from Alcatel Canada (now Thales Group). The SelTrac system is also used by the fully automated Disneyland Resort line, whereas the South Island line uses another signalling supplier, Alstom, Urbalis 400 CBTC system.
As part of RailGen 2.0 implemented from 2014 onwards to improve the standards of the rail network, the signalling systems on the older lines are to be replaced with new CBTC systems; the system used for the pre-merger network will be replaced with Alstom-Thales SelTrac whereas that for the East Rail line will be replaced by Siemens Trainguard CBTC. However, the signalling upgrades encountered a serious setback in the form of a train collision outside Central station on 18 March 2019.
Head office
The MTR Headquarters Building is located at Telford Plaza. It is a part of the larger Telford Garden complex, which was developed as part of a partnership between MTR and private development companies.
Telford Plaza held an exhibition dedicated to the history of MTR in April 2014.
Fares and tickets
After the rail merger, there are three different fare classes on the MTR: Adult, Students and Concessionary. Only children below the age of 12 and senior citizens 65 years or older are eligible for the concessionary rate on all lines. Full-time Hong Kong students between the ages of 12 and 25 qualify for the concessionary rate using a personalised Octopus Card on all lines except on Airport Express, or travel to or from cross-border stations (Lo Wu/Lok Ma Chau). Children below the age of 3 travel free (unless they exceed the height range).
The fare of MTR between any two particular stations is not calculated using a particular formula, and must be looked up from the fare table. Fares for the Airport Express Line are significantly higher. Services to checkpoint termini are also more expensive than ordinary fares, as are journeys that require a harbour crossing than are journeys that do not. Adult fares range from HK$3.6 to $52.6 (US$0.46–6.74). Concessionary fares are usually half the adult fare, and range from HK$1.50 to $27.00. Student fares are the same as child and elderly fare on the urban lines, but are the same as the Adult fares for journeys to or from checkpoint termini, and range from HK$1.50 to $51.00. The fare is subject to adjustment in June every year.
Prior to May 2009, MTR did not provide concessionary fares for the disabled. Legislators such as social welfare constituency legislator Fernando Cheung Chiu-hung and those from Hong Kong's Association for Democracy and People's Livelihood had for years demanded that such concessions be put in place. In May 2009, MTR eventually agreed to offer the disabled concessionary fares with HK$2 million sponsorship from Transport and Housing Bureau and under the condition that Legislative Council amends the Disability Discrimination Ordinance.
Single journey tickets can be purchased at vending machines while tourist passes, Octopus cards and other special tickets must be purchased at the ticket counter. Credit cards are only accepted to purchase Airport Express tickets and tourist Octopus cards from automatic vending machines located within Hong Kong airport.
There are also frequent-user passes, such as the MTR City Saver, which is valid at 67 stations and can be used for 40 trips over 40 days, and the Tuen Mun-Nam Cheong day pass, which is valid for unlimited travel for one day on a portion of the Tuen Ma line.
Octopus cards
The Octopus card is a reusable contactless stored value smart card for making electronic payments in online or offline systems in Hong Kong developed by Australian company ERG Group.
Launched in September 1997 to collect fares for the territory's mass transit system, the Octopus card system became the world's second contactless smart card system, after the Korean Upass, and has since grown into a widely used payment system for all public transport in Hong Kong. Octopus's success has led to the development of the Navigo card in Paris, the Oyster card in London, the Opal card in New South Wales, NETS FlashPay, EZ-Link in Singapore, and other similar systems.
The Octopus card has also evolved for use as payment in many retail shops in Hong Kong, including convenience stores, supermarkets, and fast-food restaurants. Other common Octopus payment applications include parking meters, car parks, petrol stations, vending machines, fee payment at public libraries and swimming pools, and more. The cards are also used for non-payment purposes, such as school attendance and access control for office buildings and housing estates.
Tourist pass
The Tourist Day Pass gives tourists unlimited MTR rides (with the exception of MTR Bus routes, the First Class section of the East Rail line, the Airport Express, as well as journeys to and from Lo Wu, Lok Ma Chau and Racecourse stations) for 24 hours from the point of first entry . Each pass costs HK$65 and is available at all the MTR Customer Service Centres. Tourist Day Pass must be used within 30 days upon the day of issue. The Airport Express Tourist Octopus Cards are also available. Cardholders may enjoy three days of unlimited rides on the MTR (except Airport Express, East Rail line First Class, and journeys involving Lo Wu and Lok Ma Chau stations) refundable deposit of HK$50 and choice of either a single (HK$220) or round trip (HK$300) on the Airport Express.
Other fares
A touchless smart card system is used for single journey tickets. These tickets are pre-paid for between pre-determined stations, and are good for only one trip. There are no return tickets, except on the Airport Express. As of mid-2013, less than five per cent of MTR customers travelled on single journey tickets.
Fares for the Airport Express are substantially different from main line fares. Apart from single tickets, same-day return tickets (same price as a single), and one-month return tickets are also available.
A one-day pass was able to be purchased for unlimited travel to and from Hong Kong Disneyland within the same day, from 2005 to 2011, and cost HK$50. This pass could be purchased from any MTR Customer Service Centres or Airport Express Customer Service Centres.
Ticket recommendation
Ticket Suggestion and Route Suggestion functions are available on the MTR website; based on trip destination and travel pattern, they can recommend the lowest price ticket type for daily and non-daily commuters.
Third parties, such as MTR Service Update, have also developed ticket recommendation capabilities, claiming to be more user-friendly and fare-saving. The Checkfare function at MTR Service Update can recommend whether to interchange at Tsim Sha Tsui or East Tsim Sha Tsui, to receive a better discount.
Performance
Since the merger in 2007, MTR has consistently achieved a 99.9% on-time rate, meaning 999 of every 1,000 passengers arrives at their destination within 5 minutes of scheduled time. In 2013, out of the 5.2 million passengers the MTR averaged each workday, 5.195 million passengers were considered to have arrived "on time". This makes the MTR one of the most efficient major public transport networks on the planet. MTR must report all delays of more than eight minutes to the government. There were 143 reportable incidents in 2013. MTR is fined HK$1 million for having delays of 30 minutes to an hour, with higher fines for longer delays.
Regulations and safety
According to the Mass Transit Railway By-laws, eating, drinking, or smoking are not allowed in the paid area of stations or in trains. Offenders will be fined up to HK$5000.
Various campaigns and activities are taken to help ensure that the MTR is a safe system to travel on. Poster campaigns displaying information on topics such as escalator safety are a common sight in all MTR stations, and announcements are made regularly as safety reminders to travelling passengers. By-laws were also introduced to deter potentially dangerous actions on the MTR, such as the ban on flammable goods on the MTR and rushing into trains when the doors are closing. Penalties ranging from fines to imprisonment have been imposed for such offences.
Police officers patrol the trains and stations, and police posts are available at some stations. The Hong Kong Police Force has a Railway District responsible for the MTR. Closed-circuit television cameras are installed in stations and on some of the newer trains.
The entire Tung Chung line and Airport Express, as well as the stations added by the Tseung Kwan O line, has platform screen doors (PSDs), ordered from Swiss glass door manufacturer Kaba Gilgen AG, installed upon construction. So does the entire Tuen Ma line, inherited from KCR. These doors make platforms safer by preventing people from falling onto the rails, even though MTRCL did not heavily promote it directly. However, the primary motivation was to separate the stations from the tunnels, hence allowing substantial energy savings on station air-conditioning and tunnel ventilation. Automatic platform gates (APGs) have also been installed at the Sunny Bay and Disneyland Resort stations. Their heights are half of the PSDs and only prevent people from falling onto the rails. MTR has finished installing the APGs on all of the above-ground stations of the MTR except on the East Rail line; they will be installed there as part of the Sha Tin to Central Link project.
In June 2000, MTRCL proceeded with plans to retrofit 2,960 pairs of platform screen doors at all 30 underground stations on the Kwun Tong, Tsuen Wan, and Island lines in a six-year programme. The programme made MTR the world's first railway to undertake the retrofitting of PSDs on a passenger-carrying system already in operation. A prototype design was first introduced at Choi Hung station in the 3rd quarter of 2001. The scheme was completed in October 2005, ahead of the forecast completion date in 2006. MTRCL said that part of the cost had to be assumed by passengers. HK$0.10 per passenger trip was levied on Octopus card users to help fund the HK$2 billion retrofit programme. This levy was ended in 2013 after raising more than HK$1 billion.
Visual identity
The MTR visual identity, which includes logo, vehicle livery, signage, route maps and passenger information, was updated in 1995–1998 by Lloyd Northover, the British design consultancy founded by John Lloyd and Jim Northover.
MTR has a mascot named "Kee Gor", modeled after a former MTR employee who headed operations.
Social outreach
Art promotion
With the objective "not only bring MTR passengers more time for life, but also more time for art", the Art in MTR Initiative has been a success since its reception in 1998, where the Airport Express Artwork Programme was the pioneer project. Thereafter, live performances, art exhibitions, display of artwork by established and emerging artists, students and young children have been brought into the MTR stations. MTRCL have even made art part of the station architecture when building new stations or renovating existing ones. Artworks are exhibited in different forms on the network, including "arttube", open art gallery, community art galleries, roving art, living art, and art in station architecture.
MTR Hong Kong Race Walking
MTR and Hong Kong Association of Athletics Affiliates have jointly hosted MTR Hong Kong Race Walking annually in spring since 2005. The race walking competition aims at promoting healthy living in Hong Kong. The race begins and ends on the ground above Central MTR station, namely Chater Garden, Chater Road, Ice House Street and Des Voeux Road Central in Central. There is a fun walk apart from the regular competition. The event attracted over 800 participants in 2005 and 1,500 in 2012. The event is attended not only by Hongkongers, but also athletes from various countries. The race raises fund for Better Health for a Better Hong Kong, a Hospital Authority project for the working population.
Controversies
Tree removal
The MTR Corporation came under fire in June 2011 after their work on the cross-border high-speed railway line encroached on a conservation area in Pat Heung, Yuen Long. 34 trees were felled and an entire slope was concreted over in the conservation area. The Environmental Protection Department issued summonses to the corporation for offences under the Environmental Impact Assessment Ordinance. In September 2011, a fine of HK$15,000 was imposed by the court. The MTR Corporation admitted that 34 trees were felled by mistake; all were common native woodland species and no rare tree species were affected. The corporation said owing to a technical misalignment of relevant drawings, the plan submitted to the Environmental Protection Department did not include the part of the Conservation Area which was included in the gazettal plan of their works. The corporation became aware that part of the approved tree removal works may have encroached onto the Conservation Area during construction, and proactively reported the situation to the government. Evaluation and measures have been taken to prevent similar incidents from happening again.
The MTR Corporation came under fire again in September 2011 after felling dozens of trees in Admiralty as part of construction work for the South Island line. Green activists denounced the tree felling as "unprofessional", and Ken So Kwok-yin, chief executive of the Conservancy Association and a certified tree arborist, said that the explanations offered by the MTR Corporation as to why the trees were felled were "unacceptable". The MTR Corporation is felling approximately 4,000 trees for the construction of the South Island line, raising concerns from environmental groups and the public about its commitment to protecting Hong Kong's natural environment.
Limits on oversized luggage
The corporation has limits on the size of items allowed on trains. The MTR system is facing pressure from increasing numbers of parallel traders who carry oversized baggage onto trains for resale in China. The corporation has been criticised for allowing parallel traders to board trains with massive bags, causing undue congestion and inconvenience to residents of the North District.
Furthermore, the corporation accused of double standards in enforcement when images of cross-border smugglers pushing overladen trollies appeared on social network sites on a regular basis, whilst local students carrying large musical instruments were reported to have been stopped and issued with written warnings. Leading musicians joined in the criticism of MTR's stance on large instruments; some citizens invited players of cellos and other large instruments to congregate on 3 October 2015 with their equipment at Tai Wai station, where the majority of these instances occurred.
Following the public uproar, MTR issued a press release in the early hours acknowledging discontent and announcing a one-month review of the policy on oversized items to see whether there was room for fine-tuning that would not compromise on passenger safety. The corporation said that staff would continue executing existing policy until any revisions are made.
October 2018 disruption
On 16 October 2018, four MTR lines suffered delays simultaneously, an unprecedented disruption to railway services. MTR stated that initial investigations showed that the problems were related to the computers that control the signalling system, and an in-depth investigation would be carried out.
Cathay Pacific advertisement
In May 2019, the MTR Corporation and the Airport Authority Hong Kong reportedly refused to display a Cathay Pacific advertisement featuring two men holding hands due to its LGBT message.
Involvement in 2019–20 Hong Kong protests
Yuen Long attack
On 21 July 2019, a mob of men dressed in white and carrying wooden sticks and metal pipes entered the MTR's Yuen Long station and assaulted people indiscriminately. The attack is largely believed to have been carried out by pro-Beijing paid thugs. One pregnant woman was hurt and found lying on the floor, and journalists were also attacked. The mob entered the paid area and attacked commuters aboard a train, which was unable to depart. Over 40 people were sent to hospital. After the incident, pro-Beijing legislator Junius Ho was accused of supporting the attack.
Prince Edward station attack
On 31 August 2019, during the anti-extradition bill protests, Special Tactical Squad officers of the Hong Kong Police Force entered Prince Edward station and attacked people inside. They fired tear gas inside the station and trains, violating guidelines on the use of such products in enclosed spaces. Bystanders were caught in the operation and it has generally been deemed a brutal attempt to stop the protests. Widespread rumours of civilian deaths at the station circulated after discrepancies were noted regarding the number of injuries. The MTR refused to provide CCTV footage filmed during the incident, helping to perpetuate these rumours.
Halting of services
On multiple occasions during the 2019–20 Hong Kong protests, MTR has sealed off stations close to locations of protest before their starting time. Supporters of the protests have thus criticised MTR of intentionally impeding the public from attending protests and unnecessarily affecting civilians, giving MTR the nickname "CCP railway" (Chinese: 黨鐵; ).
Following a clash between police and protestors in the Yuen Long station on 21 August 2019, the Chinese state media People's Daily published a commentary accusing MTR of "conspiring with protestors" by "arranging special free trains for rioters to escape". In response to the accusation, MTR issued a statement declaring that it will close stations under emergency situations in the future. After that, MTR has on multiple occasions closed off stations close to ongoing protests, for example closing the Lam Tin station, Kwun Tong station and Ngau Tau Kok station on 24 August 2020 before the starting time of a permitted demonstration in Kwun Tong. Similar incidents of varying scale have occurred multiple times later.
Arrangement of train for riot police
On 24 August 2019, MTR arranged a special train exclusively to carry riot police to Kowloon Bay station, which was closed to the public at that time due to the demonstration nearby at Kwun Tong. This has led to criticisms that MTR is assisting the government in oppressing the freedom of assembly and the freedom of expression in Hong Kong.
| Technology | China | null |
267247 | https://en.wikipedia.org/wiki/River%20dolphin | River dolphin | River dolphins are a polyphyletic group of fully aquatic mammals that reside exclusively in freshwater or brackish water. They are an informal grouping of dolphins, which itself is a paraphyletic group within the infraorder Cetacea. Extant river dolphins are placed in two superfamilies, Platanistoidea and Inioidea. They comprise the families Platanistidae (the South Asian dolphins), the possibly extinct Lipotidae (Yangtze River dolphin), Iniidae (the Amazonian dolphins) and Pontoporiidae. There are five extant species of river dolphins. River dolphins, alongside other cetaceans, belong to the clade Artiodactyla, with even-toed ungulates, and their closest living relatives the hippopotamuses, from which they diverged about 40 million years ago. Specific types of dolphins can be pink.
River dolphins are relatively small compared to other dolphins, having evolved to survive in warm, shallow water and strong river currents. They range in size from the long South Asian river dolphin to the and Amazon river dolphin. Several species exhibit sexual dimorphism, in that the females are larger than the males. They have streamlined bodies and two limbs that are modified into flippers. River dolphins use their conical-shaped teeth and long beaks to capture fast-moving prey in murky water. They have well-developed hearing that is adapted for both air and water; they do not really rely on vision since the water they swim in is usually very muddy. Instead, they tend to rely on echolocation when hunting and navigating. These species are well-adapted to living in warm, shallow waters, and, unlike other cetaceans, have little to no blubber.
River dolphins are not very widely distributed; they are all restricted to certain rivers or deltas. This makes them extremely vulnerable to habitat destruction. River dolphins feed primarily on fish. Male river dolphins typically mate with multiple females every year, but females only mate every two to three years. Calves are typically born in the spring and summer months and females bear all the responsibility for raising them. River dolphins produce a variety of vocalizations, usually in the form of clicks and whistles.
River dolphins are rarely kept in captivity; breeding success has been poor and the animals often die within a few months of capture. , there was only one river dolphin in captivity.
Taxonomy and evolution
Classification
Four families of river dolphins (Iniidae, Pontoporiidae, Lipotidae and Platanistidae) are currently recognized, comprising three superfamilies (Inioidea, Lipotoidea and Platanistoidea). Platanistidae, containing the two subspecies of South Asian river dolphin, is the only living family in the superfamily Platanistoidea. Previously, many taxonomists had assigned all river dolphins to a single family, Platanistidae, and treated the Ganges and Indus river dolphins as separate species. A December 2006 survey found no members of Lipotes vexillifer (commonly known as the baiji, or Chinese river dolphin) and declared the species functionally extinct.
The current classification of river dolphins is as follows:
Superfamily Platanistoidea
Family Platanistidae
Genus Platanista
South Asian river dolphin, Platanista gangetica, with two subspecies
Ganges river dolphin (susu), P. g. gangetica
Indus river dolphin (bhulan), P. g. minor
Family †Allodelphinidae (Oligocene - Miocene)
Family †Squalodelphinidae (Oligocene to Miocene)
Family †Squalodontidae (Oligocene to Miocene)
Family †Waipatiidae (Oligocene to Miocene)
Superfamily Inioidea
Family Iniidae
Genus Inia
Amazon river dolphin (boto), Inia geoffrensis
Inia geoffrensis geoffrensis
Inia geoffrensis humbotiana
Araguaian river dolphin, Inia araguaiaensis
Bolivian river dolphin, Inia boliviensis
Genus †Meherrinia (late Miocene)
Genus †Isthminia (Miocene)
Family Pontoporiidae
Genus †Auroracetus
†Auroracetus bakerae
Genus Pontoporia
La Plata dolphin (Franciscana), Pontoporia blainvillei
Superfamily Lipotoidea
Family Lipotidae
Genus Lipotes
Baiji (or Chinese river dolphin), Lipotes vexillifer
In 2012 the Society for Marine Mammalogy began considering the Bolivian (Inia geoffrensis boliviensis) and Amazonian (Inia geoffrensis geoffrensis) subspecies as full species Inia boliviensis and Inia geoffrensis, respectively; however, much of the scientific community, including the IUCN, continue to consider the Bolivian population to be a subspecies of Inia geoffrensis.
In October 2014, the Society for Marine Mammalogy took Inia boliviensis and Inia araguaiaensis off their list of aquatic mammal species and subspecies and currently does not recognize these species-level separations.
Evolution
River dolphins are members of the infraorder Cetacea, which are descendants of land-dwelling mammals of the order Artiodactyla (even-toed ungulates). They are related to the Indohyus, an extinct chevrotain-like ungulate, from which they split approximately 48 million years ago.
The primitive cetaceans, or archaeocetes, first took to the sea approximately 49 million years ago and became fully aquatic by 5–10 million years later. It is unknown when river dolphins first ventured back into fresh water.
River dolphins are thought to have relictual distributions, that is, their ancestors originally occupied marine habitats, but were then displaced from these habitats by modern dolphin lineages. Many of the morphological similarities and adaptations to freshwater habitats arose due to convergent evolution; thus, a grouping of all river dolphins is polyphyletic. Amazon river dolphins are actually more closely related to oceanic dolphins than to South Asian river dolphins. Isthminia panamensis is an extinct genus and species of river dolphin, living 5.8 to 6.1 million years ago. Its fossils were discovered near Piña, Panama.
River dolphin has been considered a taxonomic description, suggesting an evolutionary relationship among the group, although it is now known that they form two distinct clades. 'True' river dolphins are descendants of ancient evolutionary lineages that evolved in freshwater environments.
Some species of cetacean live in rivers and lakes, but are more closely related to oceanic dolphins or porpoises and entered fresh water more recently. Such species are considered facultative freshwater cetaceans as they can use both marine and freshwater environments. These include species such as the Irrawaddy dolphin, Orcaella brevirostris, found in the Mekong, Mahakam, the Irrawaddy Rivers, as well as the Yangtze finless porpoise Neophocaena phocaenoides asiaeorientalis. Some oceanic cetacean populations are known to live semi-permanently in river and estuarine systems, such as the Indo-Pacific bottlenose dolphin group resident in the Swan River of Western Australia which travel as far inland as Belmont.
The tucuxi (Sotalia fluviatilis) in the Amazon River is another species descended from oceanic dolphins; however, it does not perfectly fit the label of 'facultative' either, as it occurs only in fresh water. The tucuxi was until recently considered conspecific with the Guiana dolphin (Sotalia guianensis), which inhabits marine waters. It may also be true for the Irrawaddy dolphin and the finless porpoise that, although the species may be found in both freshwater and marine environments, individual animals found in rivers may not be able to survive in the ocean, and vice versa. The tucuxi is currently classified as an oceanic dolphin (Delphinidae).
The Franciscana (Pontoporia blainvillei) has shown a converse evolutionary pattern, and has an ancient evolutionary lineage in freshwater, but inhabits estuarine and coastal waters.
Biology
Anatomy
River dolphins have a torpedo shaped body with a flexible neck, limbs modified into flippers, non-existent external ear flaps, a tail fin, and a small bulbous head. River dolphin skulls have small eye orbits, a long snout and eyes placed on the sides of the head. River dolphins are rather small, ranging in size from the long South Asian river dolphin to the and Amazon river dolphin. They all have female-biased sexual dimorphism apart from Amazon river dolphin, with the females being larger than the males. River dolphins are polygynous, meaning male river dolphins typically mate with multiple females every year, but females only mate every two to three years. Calves are typically born in the spring and summer months and females bear all the responsibility for raising them.
River dolphins have conical teeth, used to catch swift prey such as small river fish. They also have very long snouts, with some measuring , four times longer than most of their oceanic counterparts. They have a two-chambered stomach that is similar in structure to that of terrestrial carnivores. They have fundic and pyloric chambers. Breathing involves expelling stale air from their blowhole, followed by inhaling fresh air into their lungs. They do not have the iconic spout, as this only forms when the warm air exhaled from the lungs meets cold external air, which does not occur in their tropical habitats.
River dolphins have a relatively thin layer of blubber. Blubber can help with buoyancy, protection from predators (they would have a hard time getting through a thick layer of fat), energy for leaner times, and insulation from harsh climates. The habitats of river dolphins lack these needs.
Locomotion
River dolphins have two flippers and a tail fin. These flippers contain four digits. Although river dolphins do not possess fully developed hind limbs, some possess discrete rudimentary appendages, which may contain feet and digits. River dolphins are slow swimmers in comparison to oceanic dolphins, which can travel at speeds up to ; the tucuxi can only travel at about . Unlike other cetaceans, their neck vertebrae are not fused together, meaning they have greater flexibility than other non-terrestrial aquatic mammals, at the expense of speed. This means they can turn their head without actually moving their entire body. When swimming, river dolphins rely on their tail fins to propel themselves through the water. Flipper movement is continuous. River dolphins swim by moving their tail fins and lower bodies up and down, propelling themselves through vertical movement, while their flippers are mainly used for steering. All species have a dorsal fin.
Senses
The ears of river dolphins have specific adaptations to their aquatic environment. In humans, the middle ear works as an impedance equalizer between the outside air's low impedance and the cochlear fluid's high impedance. In river dolphins, and other cetaceans, there is no great difference between the outer and inner environments. Instead of sound passing through the outer ear to the middle ear, river dolphins receive sound through the throat, from which it passes through a low-impedance fat-filled cavity to the inner ear. The ear is acoustically isolated from the skull by air-filled sinus pockets, which allows for greater directional hearing underwater. Dolphins send out high frequency clicks from an organ known as a melon. This melon consists of fat, and the skull of any such creature containing a melon will have a large depression. This allows river dolphins to produce biosonar for orientation. They are so dependent on echolocation that they can survive even if they are blind. Beyond locating an object, echolocation also provides the animal with an idea on the object's shape and size, though how exactly this works is not yet understood. The small hairs on the rostrum of the Amazon river dolphin are believed to function as a tactile sense, possibly to compensate for their poor eyesight.
River dolphins have very small eyes for their size, and do not have a very good sense of sight. In addition, the eyes are placed on the sides of the head, so the vision consists of two fields, rather than a binocular view like humans have. When river dolphins surface, their lens and cornea correct the nearsightedness that results from the refraction of light. They have both rod and cone cells, meaning they can see in both dim and bright light. Most river dolphins have slightly flattened eyeballs, enlarged pupils (which shrink as they surface to prevent damage), slightly flattened corneas and a tapetum lucidum; these adaptations allow for large amounts of light to pass through the eye and, therefore, a very clear image of the surrounding area. They also have glands on their eyelids and an outer corneal layer that act as protection for the cornea.
Olfactory lobes are absent in river dolphins, suggesting that they have no sense of smell.
River dolphins are not thought to have a sense of taste, as their taste buds are atrophied or missing altogether. However, some dolphins have preferences between different kinds of fish, indicating some sort of attachment to taste.
Interactions with humans
Threats
Development
Development and agriculture have had devastating impacts on the habitats on river dolphins. The total population of Araguaian river dolphins is estimated to be between 600 and 1,500 individuals, and genetic diversity is limited. The ecology of their habitat has been adversely affected by agricultural, ranching and industrial activities, as well as by the use of dams for hydroelectric power. The inhabited section of the Araguaia River probably extends over about out of a total length of . The Tocantins river habitat is fragmented by six hydroelectric dams, so the population there is at particular risk. Its probable eventual IUCN status is vulnerable or worse.
Both subspecies of South Asian river dolphins have been very adversely affected by human use of the river systems in the subcontinent. Irrigation has lowered water levels throughout both subspecies' ranges. Poisoning of the water supply from industrial and agricultural chemicals may have also contributed to population decline. Perhaps the most significant issue is the building of more than 50 dams along many rivers, causing the segregation of populations and a narrowed gene pool in which the dolphins can breed. Currently, three subpopulations of Indus river dolphins are considered capable of long-term survival if protected.
As China developed economically, pressure on the baiji river dolphin grew significantly. Industrial and residential waste flowed into the Yangtze. The riverbed was dredged and reinforced with concrete in many locations. Ship traffic multiplied, boats grew in size, and fishermen employed wider and more lethal nets. Noise pollution caused the nearly blind animal to collide with propellers. Stocks of the dolphin's prey declined drastically in the late 20th century, with some fish populations declining to one thousandth of their pre-industrial levels. In the 1950s, the population was estimated at 6,000 animals, but declined rapidly over the subsequent five decades. Only a few hundred were left by 1970. Then the number dropped down to 400 by the 1980s and then to 13 in 1997 when a full-fledged search was conducted. On December 13, 2006, the baiji (Lipotes vexillifer) was declared "functionally extinct", after a 45-day search by leading experts in the field failed to find a single specimen. The last verified and widely accepted sighting was in September 2004, but it has been allegedly seen and photographed by Chinese citizens on four occasions since then.
Competition
The region of the Amazon in Brazil has an extension of containing diverse fundamental ecosystems. One of these ecosystems is a floodplain, or a várzea forest, and is home to a large number of fish species which are an essential resource for human consumption. The várzea is also a major source of income through excessive local commercialized fishing. Várzea consist of muddy river waters containing a vast number and diversity of nutrient-rich species. The abundance of distinct fish species lures the Amazon River dolphin into the várzea areas of high water occurrences during the seasonal flooding.
In addition to attracting predators such as the Amazon river dolphin, these high-water occurrences are an ideal location to draw in the local fisheries. Human fishing activities directly compete with the dolphins for the same fish species, the tambaqui (Colossoma macropomum) and the pirapitinga (Piaractus brachypomus), resulting in deliberate or unintentional catches of the Amazon river dolphin. The local fishermen overfish, and when the Amazon river dolphins remove the commercialized fish from the nets and lines, it damages the equipment and the capture and causes a negative reaction from the local fishermen.
The Brazilian Institute of Environment and Renewable Natural Resources prohibit fishermen from killing the Amazon river dolphin, yet they are not compensated for the damage to their equipment and the loss of their catch.
Bycatch
During the process of catching the commercialized fish, the Amazon river dolphins get caught in the nets and exhaust themselves until they die, or the local fishermen deliberately kill the dolphins that become entangled in their nets. The carcasses are discarded, consumed, or used as bait to attract a scavenger catfish, the piracatinga (Calophysus macropterus). The use of the Amazon river dolphin carcass as bait for the piracatinga dates back from 2000. The increasing consumption demand by the local inhabitants and Colombia for the piracatinga has created a market for distribution of the Amazon river dolphin carcasses to be used as bait throughout these regions.
For example, of the 15 dolphin carcasses found in the Japurá River in 2010–2011 surveys, 73% of the dolphins were killed for bait, disposed of, or abandoned in entangled gillnets. The data does not fully represent the actual overall number of deaths of the Amazon river dolphins, whether accidental or intentional, because a variety of factors make it extremely complicated to record and medically examine all the carcasses. Scavenger species feed upon them and the complexity of the river currents makes it nearly impossible to locate all the carcasses. More importantly, the local fishermen do not report these deaths out of fear that legal action will be taken against them, as the Amazon river dolphin and other cetaceans are protected under the Brazilian federal law, prohibiting any takes, harassments, and kills of the species.
Conservation
The Global Declaration for River Dolphins was signed by nine countries on 24 October 2023, a date chosen as it is known as the International River Dolphin Day. This pact is intended to promote research and cooperation between countries with river dolphin populations. It is hoped that five further countries will join.
In captivity
A baiji conservation dolphinarium was established at the Institute of Hydrobiology (IHB) in Wuhan in 1992. This was planned as a backup to any other conservation efforts by producing an area completely protected from any threats, and where the baiji could be easily observed. The site includes an indoor and outdoor holding pool, a water filtration system, food storage and preparation facilities, research labs and a small museum. The aim is to also generate income from tourism which can be put towards the baiji plight. The pools are not very large, only kidney shaped tanks with dimensions of arc width and depth, diameter, deep and diameter, deep, and are not capable of holding many baijis at one time.
Douglas Adams and Mark Carwardine documented their encounters with the endangered animals on their conservation travels for the BBC programme Last Chance to See. The book by the same name, published in 1990, included pictures of a captive specimen, a male named Qi Qi (淇淇) that lived in the Wuhan Institute of Hydrobiology dolphinarium from 1980 to July 14, 2002. Discovered by a fisherman in Dongting Lake, he became the sole resident of the Baiji Dolphinarium (白鱀豚水族馆) beside East Lake. A sexually mature female was captured in late 1995, but died after half a year in 1996 when the Shishou Tian-e-Zhou Baiji Semi-natural Reserve (石首半自然白鱀豚保护区), which had contained only finless porpoises since 1990, was flooded.
The Amazon river dolphin has historically been kept in dolphinariums. Today, only three exist in captivity: one in Acuario de Valencia in Venezuela, one in Zoologico de Guistochoca in Peru, and one in Duisburg Zoo in Germany. Several hundred were captured between the 1950s and 1970s, and were distributed in dolphinariums throughout the US, Europe, and Japan. Around 100 went to US dolphinariums, and of that, only 20 survived; the last (named Chuckles) died in Pittsburgh Zoo in 2002.
In mythology
Asia
In Hindu mythology, the Ganges river dolphin is associated with Ganga, the deity of the Ganges river. The dolphin is said to be one of the creatures which heralded the goddess' descent from the heavens, and Ganga's mount, the Makara, is sometimes depicted as a dolphin.
In Chinese mythology, the baiji has many origin stories.
For example, near the mouth of the Yangtze, the baiji was a princess that had lost her parents and had lived with her stepfather, whom she had longed to get away from. The stepfather wanted to trade her since she would be sold for a great sum of money, but as they were crossing the river to get to the trader, a storm rolled in. The enraged stepfather tried to take her, but she plunged herself into the river, was transformed into a dolphin before she drowned, and swam away from her abusive stepfather, who also fell in and was transformed into a porpoise.
In another story, the baiji was the daughter of a general deported from the city of Wuhan during a war who ran away while her father was in duty. Later, the general met a woman who told him how her father was a general. When he realized that she was his daughter, he threw himself into the river out of shame, and his daughter ran after him and also fell into the river. Before they were drowned, the daughter was transformed into a dolphin and the general into a porpoise.
South America
Amazon river dolphins, known by the natives as the boto, encantados or toninhas, are very prevalent in the mythology of the native South Americans. They are often characterized in mythology with superior musical ability, seductiveness and love of sex that often results in illegitimate children, and attraction to parties. Despite the fact that the Encante are said to come from a utopia full of wealth and without pain or death, the encantados crave the pleasures and hardships of human societies.
Transformation into human form is said to be rare, and usually occurs at night. The encantado will often be seen running from a festa, despite protests from the others for it to stay, and can be seen by pursuers as it hurries to the river and reverts to dolphin form. When it is under human form, it wears a hat to hide its blowhole, which does not disappear with the shapeshift.
Besides the ability to shapeshift into human form, encantados frequently wield other magical abilities, such as controlling storms, hypnotizing humans into doing their will, transforming humans into encantados, and inflicting illness, insanity, and even death. Shamans often intervene in these situations.
Kidnapping is also a common theme in such folklore. Encantados are said to be fond of abducting humans with whom they fall in love, children born of their illicit love affairs, or just about anyone near the river who can keep them company, and taking them back to the Encante. The fear of this is so great among people who live near the Amazon River that both children and adults are terrified of going near the water between dusk and dawn, or entering water alone. Some who supposedly have encountered encantados while out in their canoes have been said to have gone insane, but the creatures seem to have done little more than follow their boats and nudge them from time to time.
| Biology and health sciences | Toothed whale | Animals |
267362 | https://en.wikipedia.org/wiki/Astronomical%20seeing | Astronomical seeing | In astronomy, seeing is the degradation of the image of an astronomical object due to turbulence in the atmosphere of Earth that may become visible as blurring, twinkling or variable distortion. The origin of this effect is rapidly changing variations of the optical refractive index along the light path from the object to the detector.
Seeing is a major limitation to the angular resolution in astronomical observations with telescopes that would otherwise be limited through diffraction by the size of the telescope aperture.
Today, many large scientific ground-based optical telescopes include adaptive optics to overcome seeing.
The strength of seeing is often characterized by the angular diameter of the long-exposure image of a star (seeing disk) or by the Fried parameter r0. The diameter of the seeing disk is the full width at half maximum of its optical intensity. An exposure time of several tens of milliseconds can be considered long in this context. The Fried parameter describes the size of an imaginary telescope aperture for which the diffraction limited angular resolution is equal to the resolution limited by seeing. Both the size of the seeing disc and the Fried parameter depend on the optical wavelength, but it is common to specify them for 500 nanometers.
A seeing disk smaller than 0.4 arcseconds or a Fried parameter larger than 30 centimeters can be considered excellent seeing. The best conditions are typically found at high-altitude observatories on small islands, such as those at Mauna Kea or La Palma.
Effects
Astronomical seeing has several effects:
It causes the images of point sources (such as stars), which in the absence of atmospheric turbulence would be steady Airy patterns produced by diffraction, to break up into speckle patterns, which change very rapidly with time (the resulting speckled images can be processed using speckle imaging)
Long exposure images of these changing speckle patterns result in a blurred image of the point source, called a seeing disc
The brightness of stars appears to fluctuate in a process known as scintillation or twinkling
Atmospheric seeing causes the fringes in an astronomical interferometer to move rapidly
The distribution of atmospheric seeing through the atmosphere (the CN2 profile described below) causes the image quality in adaptive optics systems to degrade the further you look from the location of reference star
The effects of atmospheric seeing were indirectly responsible for the belief that there were canals on Mars. In viewing a bright object such as Mars, occasionally a still patch of air will come in front of the planet, resulting in a brief moment of clarity. Before the use of charge-coupled devices, there was no way of recording the image of the planet in the brief moment other than having the observer remember the image and draw it later. This had the effect of having the image of the planet be dependent on the observer's memory and preconceptions which led the belief that Mars had linear features.
The effects of atmospheric seeing are qualitatively similar throughout the visible and near infrared wavebands. At large telescopes the long exposure image resolution is generally slightly higher at longer wavelengths, and the timescale (t0 - see below) for the changes in the dancing speckle patterns is substantially lower.
Measures
There are three common descriptions of the astronomical seeing conditions at an observatory:
The full width at half maximum (FWHM) of the seeing disc
r0 (the size of a typical "lump" of uniform air within the turbulent atmosphere) and t0 (the time-scale over which the changes in the turbulence become significant)
The CN2 profile
These are described in the sub-sections below:
The full width at half maximum (FWHM) of the seeing disc
Without an atmosphere, a small star would have an apparent size, an "Airy disk", in a telescope image determined by diffraction and would be inversely proportional to the diameter of the telescope. However, when light enters the Earth's atmosphere, the different temperature layers and different wind speeds distort the light waves, leading to distortions in the image of a star. The effects of the atmosphere can be modeled as rotating cells of air moving turbulently. At most observatories, the turbulence is only significant on scales larger than r0 (see below—the seeing parameter r0 is 10–20 cm at visible wavelengths under the best conditions) and this limits the resolution of telescopes to be about the same as given by a space-based 10–20 cm telescope.
The distortion changes at a high rate, typically more frequently than 100 times a second. In a typical astronomical image of a star with an exposure time of seconds or even minutes, the different distortions average out as a filled disc called the "seeing disc". The diameter of the seeing disk, most often defined as the full width at half maximum (FWHM), is a measure of the astronomical seeing conditions.
It follows from this definition that seeing is always a variable quantity, different from place to place, from night to night, and even variable on a scale of minutes. Astronomers often talk about "good" nights with a low average seeing disc diameter, and "bad" nights where the seeing diameter was so high that all observations were worthless.
The FWHM of the seeing disc (or just "seeing") is usually measured in arcseconds, abbreviated with the symbol (″). A 1.0″ seeing is a good one for average astronomical sites. The seeing of an urban environment is usually much worse. Good seeing nights tend to be clear, cold nights without wind gusts. Warm air rises (convection), degrading the seeing, as do wind and clouds. At the best high-altitude mountaintop observatories, the wind brings in stable air which has not previously been in contact with the ground, sometimes providing seeing as good as 0.4".
r0 and t0
The astronomical seeing conditions at an observatory can be conveniently described by the parameters r0 and t0.
For telescopes with diameters smaller than r0, the resolution of long-exposure images is determined primarily by diffraction and the size of the Airy pattern and thus is inversely proportional to the telescope diameter.
For telescopes with diameters larger than r0, the image resolution is determined primarily by the atmosphere and is independent of telescope diameter, remaining constant at the value given by a telescope of diameter equal to r0. r0 also corresponds to the length-scale over which the turbulence becomes significant (10–20 cm at visible wavelengths at good observatories), and t0 corresponds to the time-scale over which the changes in the turbulence become significant. r0 determines the spacing of the actuators needed in an adaptive optics system, and t0 determines the correction speed required to compensate for the effects of the atmosphere.
The parameters r0 and t0 vary with the wavelength used for the astronomical imaging, allowing slightly higher resolution imaging at longer wavelengths using large telescopes.
The seeing parameter r0 is often known as the Fried parameter, named after David L. Fried. The atmospheric time constant t0 is often referred to as the Greenwood time constant, after Darryl Greenwood.
Mathematical description of r0 and t0
Mathematical models can give an accurate model of the effects of astronomical seeing on images taken through ground-based telescopes. Three simulated short-exposure images are shown at the right through three different telescope diameters (as negative images to highlight the fainter features more clearly—a common astronomical convention). The telescope diameters are quoted in terms of the Fried parameter (defined below). is a commonly used measurement of the astronomical seeing at observatories. At visible wavelengths, varies from 20 cm at the best locations to 5 cm at typical sea-level sites.
In reality, the pattern of blobs (speckles) in the images changes very rapidly, so that long-exposure photographs would just show a single large blurred blob in the center for each telescope diameter. The diameter (FWHM) of the large blurred blob in long-exposure images is called the seeing disc diameter, and is independent of the telescope diameter used (as long as adaptive optics correction is not applied).
It is first useful to give a brief overview of the basic theory of optical propagation through the atmosphere. In the standard classical theory, light is treated as an oscillation in a field . For monochromatic plane waves arriving from a distant point source with wave-vector :
where is the complex field at position and
time , with real and imaginary parts corresponding to the electric and magnetic field components, represents a phase offset,
is the frequency of the light determined by , and is the amplitude of the light.
The photon flux in this case is proportional to the square of the amplitude , and the optical phase corresponds to the complex argument of . As wavefronts pass through the Earth's atmosphere they may be perturbed by refractive index variations in the atmosphere. The diagram at the top-right of this page shows schematically a turbulent layer in the Earth's atmosphere perturbing planar wavefronts before they enter a telescope. The perturbed wavefront may be related at any given instant to the original planar wavefront in the following way:
where represents the fractional change in wavefront amplitude and
is the change in wavefront phase introduced by the atmosphere. It is important to emphasise that and describe the effect of the Earth's atmosphere, and the timescales for any changes in these functions will be set by the speed of refractive index fluctuations in the atmosphere.
The Kolmogorov model of turbulence
A description of the nature of the wavefront perturbations introduced by the atmosphere is provided by the Kolmogorov model developed by Tatarski, based partly on the studies of turbulence by the Russian mathematician Andrey Kolmogorov. This model is supported by a variety of experimental measurements and is widely used in simulations of astronomical imaging. The model assumes that the
wavefront perturbations are brought about by variations in the refractive index of the atmosphere. These refractive index variations lead directly to phase fluctuations described by , but any amplitude fluctuations are only brought about as a second-order effect while the perturbed wavefronts propagate from the perturbing atmospheric layer to the telescope. For all reasonable models of the Earth's atmosphere at optical and
infrared wavelengths the instantaneous imaging performance is dominated by the phase fluctuations . The amplitude fluctuations described by have negligible effect on the structure of the images seen in the focus of a large telescope.
For simplicity, the phase fluctuations in Tatarski's model are often assumed to have a Gaussian random distribution with the following second-order structure function:
where is the atmospherically induced variance between the phase at two parts of the wavefront separated by a distance in the aperture plane, and represents the ensemble average.
For the Gaussian random approximation, the structure function of Tatarski (1961) can be described in terms of a single parameter :
indicates the strength of the phase fluctuations as it corresponds to the diameter of a circular telescope aperture at which atmospheric phase perturbations begin to seriously limit the image resolution. Typical values for I band (900 nm wavelength) observations at good sites are 20–40 cm. also corresponds to the aperture diameter for which the variance of the wavefront phase averaged over the aperture comes approximately to unity:
This equation represents a commonly used definition for , a parameter frequently used to describe the atmospheric conditions at astronomical observatories.
can be determined from a measured CN2 profile (described below) as follows:
where the turbulence strength varies as a function of height above the telescope, and is the angular distance of the astronomical source from the zenith (from directly overhead).
If turbulent evolution is assumed to occur on slow timescales, then the timescale t0 is simply proportional to r0 divided by the mean wind speed.
The refractive index fluctuations caused by Gaussian random turbulence can be simulated using the following algorithm:
where is the optical phase error introduced by atmospheric turbulence, R (k) is a two-dimensional square array of independent random complex numbers which have a Gaussian distribution about zero and white noise spectrum, K (k) is the (real) Fourier amplitude expected from the Kolmogorov (or Von Karman) spectrum, Re[] represents taking the real part, and FT[] represents a discrete Fourier transform of the resulting two-dimensional square array (typically an FFT).
Turbulent intermittency
The assumption that the phase fluctuations in Tatarski's model have a Gaussian random distribution is usually unrealistic. In reality, turbulence exhibits intermittency.
These fluctuations in the turbulence strength can be straightforwardly simulated as follows:
where is a two-dimensional array which represents the spectrum of intermittency, with the same dimensions as , and where represents convolution. The intermittency is described in terms of fluctuations in the turbulence strength . It can be seen that the equation for the Gaussian random case above is just the special case from this equation with:
where is the Dirac delta function.
The profile
A more thorough description of the astronomical seeing at an observatory is given by producing a profile of the turbulence strength as a function of altitude, called a profile. profiles are generally performed when deciding on the type of adaptive optics system which will be needed at a particular telescope, or in deciding whether or not a particular location would be a good site for setting up a new astronomical observatory. Typically, several methods are used simultaneously for measuring the profile and then compared. Some of the most common methods include:
SCIDAR (imaging the shadow patterns in the scintillation of starlight)
LOLAS (a small-aperture variant of SCIDAR designed for low-altitude profiling)
SLODAR
MASS
MooSci (11-channel lunar scintillometer for ground level profiling)
RADAR mapping of turbulence
Balloon-borne thermometers to measure how quickly the air temperature is fluctuating with time due to turbulence
V2 Precision Data Collection Hub (PDCH) with differential temperature sensors use to measure atmospheric turbulence
There are also mathematical functions describing the profile. Some are empirical fits from measured data and others attempt to incorporate elements of theory. One common model for continental land masses is known as Hufnagel-Valley after two workers in this subject.
Mitigation
The first answer to this problem was speckle imaging, which allowed bright objects with simple morphology to be observed with diffraction-limited angular resolution. Later came space telescopes, such as NASA's Hubble Space Telescope, working outside the atmosphere and thus not having any seeing problems and allowing observations of faint targets for the first time (although with poorer resolution than speckle observations of bright sources from ground-based telescopes because of Hubble's smaller telescope diameter). The highest resolution visible and infrared images currently come from imaging optical interferometers such as the Navy Prototype Optical Interferometer or Cambridge Optical Aperture Synthesis Telescope, but those can only be used on very bright stars.
Starting in the 1990s, many telescopes have developed adaptive optics systems that partially solve the seeing problem. The best systems so far built, such as SPHERE on the ESO VLT and GPI on the Gemini telescope, achieve a Strehl ratio of 90% at a wavelength of 2.2 micrometers, but only within a very small region of the sky at a time.
A wider field of view can be obtained by using multiple deformable mirrors conjugated to several atmospheric heights and measuring the vertical structure of the turbulence, in a technique known as Multiconjugate Adaptive Optics.
Another cheaper technique, lucky imaging, has had good results on smaller telescopes. This idea dates back to pre-war naked-eye observations of moments of good seeing, which were followed by observations of the planets on cine film after World War II. The technique relies on the fact that every so often the effects of the atmosphere will be negligible, and hence by recording large numbers of images in real-time, a 'lucky' excellent image can be picked out. This happens more often when the number of r0-size patches over the telescope pupil is not too large, and the technique consequently breaks down for very large telescopes. It can nonetheless outperform adaptive optics in some cases and is accessible to amateurs. It does require very much longer observation times than adaptive optics for imaging faint targets, and is limited in its maximum resolution.
| Physical sciences | Basics | Astronomy |
267370 | https://en.wikipedia.org/wiki/Alligator | Alligator | An alligator, or colloquially gator, is a large reptile in the genus Alligator of the family Alligatoridae of the order Crocodilia. The two extant species are the American alligator (A. mississippiensis) and the Chinese alligator (A. sinensis). Additionally, several extinct species of alligator are known from fossil remains. Alligators first appeared during the late Eocene epoch about 37 million years ago.
The name "alligator" is likely an anglicized form of , the Spanish term for "the lizard", which early Spanish explorers and settlers in Florida called the alligator. Early English spellings of the name included allagarta and alagarto.
Evolution
Alligators and caimans split in North America during the early Tertiary or late Cretaceous (about 53 million to about 65 million years ago). The Chinese alligator split from the American alligator about 33 million years ago and probably descended from a lineage that crossed the Bering land bridge during the Neogene. The modern American alligator is well represented in the fossil record of the Pleistocene. The alligator's full mitochondrial genome was sequenced in the 1990s. The full genome, published in 2014, suggests that the alligator evolved much more slowly than mammals and birds.
Phylogeny
The genus Alligator belongs to the subfamily Alligatorinae, which is the sister taxon to Caimaninae (the caimans). Together, these two subfamilies form the family Alligatoridae. The cladogram below shows the phylogeny of alligators.
Species
Extant
Extinct
Alligator hailensis
Alligator mcgrewi
Alligator mefferdi
Alligator munensis
Alligator olseni
Alligator prenasalis
Alligator thomsoni
Description
An average adult American alligator's weight and length is and , but they sometimes grow to long and weigh over . The largest ever recorded, found in Louisiana, measured . The Chinese alligator is smaller, rarely exceeding in length. Additionally, it weighs considerably less, with males rarely over .
Adult alligators are black or dark olive-brown with white undersides, while juveniles have bright yellow or whitish stripes which sharply contrast against their dark hides, providing them additional camouflage amongst reeds and wetland grasses.
Alligators commonly live up to 50 years, but there have been examples of alligators living over 70. One of the oldest recorded alligator lives was that of Saturn, an American alligator who was hatched in 1936 in Mississippi and spent nearly a decade in Germany before spending the majority of his life at the Moscow Zoo, where he died at the age of 83 or 84 on 22 May 2020. Another one of the oldest lives on record is that of Muja, an American alligator who was brought as an adult specimen to the Belgrade Zoo in Serbia from Germany in 1937. Although no valid records exist about his date of birth, as of 2012, he was in his 80s and possibly the oldest alligator living in captivity.
Habitat
Alligators are native only to the United States and China.
American alligators are found in the southeast United States: all of Florida and Louisiana; the southern parts of Georgia, Alabama, and Mississippi; coastal South and North Carolina; East Texas, the southeast corner of Oklahoma, and the southern tip of Arkansas. Louisiana has the largest alligator population. The majority of American alligators inhabit Florida and Louisiana, with over a million alligators in each state. Southern Florida is the only place where both alligators and crocodiles live side by side.
American alligators live in freshwater environments, such as ponds, marshes, wetlands, rivers, lakes, and swamps, as well as in brackish water. When they construct alligator holes in the wetlands, they increase plant diversity and provide habitat for other animals during droughts. They are, therefore, considered an important species for maintaining ecological diversity in wetlands. Farther west, in Louisiana, heavy grazing by nutrias and muskrats is causing severe damage to coastal wetlands. Large alligators feed extensively on nutrias, and provide a vital ecological service by reducing nutria numbers.
The Chinese alligator currently is found in only the Yangtze River valley and parts of adjacent provinces and is extremely endangered, with only a few dozen believed to be left in the wild. Far more Chinese alligators live in zoos around the world than can be found in the wild. Rockefeller Wildlife Refuge in southern Louisiana has several in captivity in an attempt to preserve the species. Miami MetroZoo in Florida also has a breeding pair of Chinese alligators.
Behavior
Large male alligators are solitary territorial animals. Smaller alligators can often be found in large numbers close to each other. The largest of the species (both males and females) defend prime territory; smaller alligators have a higher tolerance for other alligators within a similar size class.
Alligators move on land by two forms of locomotion, referred to as "sprawl" and "high walk". The sprawl is a forward movement with the belly making contact with the ground and is used to transition to "high walk" or to slither over wet substrate into water. The high walk is an up-on-four-limbs forward motion used for overland travel with the belly well up from the ground. Alligators have also been observed to rise up and balance on their hind legs and semi-step forward as part of a forward or upward lunge. However, they can not walk on their hind legs.
Although the alligator has a heavy body and a slow metabolism, it is capable of short bursts of speed, especially in very short lunges. Alligators' main prey are smaller animals they can kill and eat with a single bite. They may kill larger prey by grabbing it and dragging it into the water to drown. Alligators consume food that cannot be eaten in one bite by allowing it to rot or by biting and then performing a "death roll", spinning or convulsing wildly until bite-sized chunks are torn off. Critical to the alligator's ability to initiate a death roll, the tail must flex to a significant angle relative to its body. An alligator with an immobilized tail cannot perform a death roll.
Most of the muscle in an alligator's jaw evolved to bite and grip prey. The muscles that close the jaws are powerful, but the muscles for opening their jaws are weak. As a result, an adult human can hold an alligator's jaws shut bare-handed. It is common to use several wraps of duct tape to prevent an adult alligator from opening its jaws when being handled or transported.
Alligators are generally timid towards humans and tend to walk or swim away if one approaches. This may encourage people to approach alligators and their nests, which can provoke the animals into attacking. In Florida, feeding wild alligators at any time is illegal. If fed, the alligators will eventually lose their fear of humans and will learn to associate humans with food.
Diet
The type of food eaten by alligators depends upon their age and size. When young, alligators eat fish, insects, snails, crustaceans, and worms. As they mature, progressively larger prey is taken, including larger fish such as gar, turtles, and various mammals, particularly nutrias and muskrats, as well as birds, deer, and other reptiles. Their stomachs also often contain gizzard stones. They will even consume carrion if they are sufficiently hungry. In some cases, larger alligators are known to ambush dogs, Florida panthers and black bears, making them the apex predator throughout their distribution. In this role as a top predator, it may determine the abundance of prey species, including turtles and nutrias. As humans encroach into their habitat, attacks are few but not unknown. Alligators, unlike the large crocodiles, do not immediately regard a human upon encounter as prey, but may still attack in self-defense if provoked.
Reproduction
Alligators generally mature at a length of . The mating season is in late spring. In April and May, alligators form so-called "bellowing choruses". Large groups of animals bellow together for a few minutes a few times a day, usually one to three hours after sunrise. The bellows of male American alligators are accompanied by powerful blasts of infrasound. Another form of male display is a loud head-slap. In 2010, on spring nights alligators were found to gather in large numbers for group courtship, the so-called "alligator dances".
In summer, the female builds a nest of vegetation where the decomposition of the vegetation provides the heat needed to incubate the eggs. The sex of the offspring is determined by the temperature in the nest and is fixed within seven to 21 days of the start of incubation. Incubation temperatures of or lower produce a clutch of females; those of or higher produce entirely males. Nests constructed on leaves are hotter than those constructed on wet marsh, so the former tend to produce males and the latter, females. The baby alligator's egg tooth helps it get out of its egg during hatching time. The natural sex ratio at hatching is five females to one male. Females hatched from eggs incubated at weigh significantly more than males hatched from eggs incubated at . The mother defends the nest from predators and assists the hatchlings to water. She will provide protection to the young for about a year if they remain in the area. Adult alligators regularly cannibalize younger individuals, though estimates of the rate of cannibalism vary widely. In the past, immediately following the outlawing of alligator hunting, populations rebounded quickly due to the suppressed number of adults preying upon juveniles, increasing survival among the young alligators.
Anatomy
Alligators, much like birds, have been shown to exhibit unidirectional movement of air through their lungs. Most other amniotes are believed to exhibit bidirectional, or tidal breathing. For a tidal breathing animal, such as a mammal, air flows into and out of the lungs through branching bronchi which terminate in small dead-end chambers called alveoli. As the alveoli represent dead-ends to flow, the inspired air must move back out the same way it came in. In contrast, air in alligator lungs makes a circuit, moving in only one direction through the parabronchi. The air first enters the outer branch, moves through the parabronchi, and exits the lung through the inner branch. Oxygen exchange takes place in extensive vasculature around the parabronchi.
The alligator has a similar digestive system to that of the crocodile, with minor differences in morphology and enzyme activity. Alligators have a two-part stomach, with the first smaller portion containing gastroliths. It is believed this portion of the stomach serves a similar function as it does in the gizzard of some species of birds, to aid in digestion. The gastroliths work to grind up the meal as alligators will take large bites or swallow smaller prey whole. This process makes digestion and nutrient absorption easier once the food reaches the second portion of the stomach. Once an alligator's meal has been processed it will move on to the second portion of the stomach which is highly acidic. The acidity of the stomach has been observed to increase once digestion begins. This is due to the increase in CO2 concentration of the blood, resulting from the right to left shunting of the alligators heart. The right to left shunt of the heart in alligators means the circulatory system will recirculate blood through the body instead of back to the lungs. The re-circulation of blood leads to higher CO2 concentration as well as lower oxygen affinity. There is evidence to suggest that there is increased blood flow diverted to the stomach during digestion to facilitate an increase in CO2 concentration which aids in increasing gastric acid secretions during digestion. The alligator's metabolism will also increase after a meal by up to four times its basal metabolic rate. Alligators also have highly folded mucosa in the lining of the intestines to further aid in the absorption of nutrients. The folds result in greater surface area for the nutrients to be absorbed through.
Alligators also have complex microbiomes that are not fully understood yet, but can be attributed to both benefits and costs to the animal. These microorganisms can be found in the high surface area of the mucosa folds of the intestines, as well as throughout the digestive tract. Benefits include better total health and stronger immune system. However alligators are still vulnerable to microbial infections despite the immune boost from other microbiota.
During brumation the process of digestion experiences changes due to the fasting most alligators experience during these periods of inactivity. Alligators that go long enough without a meal during brumation will begin a process called autophagy, where the animal begins to consume its fat reserves to maintain its body weight until it can acquire a sufficient meal. There is also fluctuation in the level of bacterial taxa populations in the alligator's microbial community between seasons which helps the alligator cope with different rates of feeding and activity.
Like other crocodilians, alligators have an armor of bony scutes. The dermal bones are highly vascularised and aid in calcium balance, both to neutralize acids while the animal cannot breathe underwater and to provide calcium for eggshell formation.
Alligators have muscular, flat tails that propel them while swimming.
The two kinds of white alligators are albino and leucistic. These alligators are practically impossible to find in the wild. They could survive only in captivity and are few in number. The Aquarium of the Americas in New Orleans has leucistic alligators found in a Louisiana swamp in 1987.
Human uses
Alligators are raised commercially for their meat and their skin, which when tanned is used for the manufacture of luggage, handbags, shoes, belts, and other leather items. Alligators also provide economic benefits through the ecotourism industry. Visitors may take swamp tours, in which alligators are a feature. Their most important economic benefit to humans may be the control of nutrias and muskrats.
Alligator meat is also consumed by humans.
Differences from crocodiles
While there are rules of thumb for distinguishing alligators from crocodiles, all of them admit exceptions. Such general rules include:
Exposed vs. interdigitated teeth: The easiest way to distinguish crocodiles from alligators is by looking at their jaw line. The teeth on the lower jaw of an alligator fit into sockets in the upper jaw, leaving only the upper teeth visible when the mouth is closed. The teeth on the lower jaw of a crocodile fit into grooves on the outside of the top jaw, making both the upper and lower teeth visible when the mouth is closed, thus creating a "toothy grin."
Shape of the nose and jaw: Alligators have wider, shovel-like, U-shaped snouts, while crocodile snouts are typically more pointed or V-shaped. The alligators' broader snouts have been contentiously thought to allow their jaws to withstand the stress of cracking open the shells of turtles and other hard-shelled animals that are widespread in their environments. A 2012 study found very little correlation between bite force and snout shape amongst 23 tested crocodilian species.
Functioning salt glands: Crocodilians have modified salivary glands called salt glands on their tongues, but while these organs still excrete salt in crocodiles and gharials, those in most alligators and caimans have lost this ability, or excrete it in only extremely small quantities. The ability to excrete excess salt allows crocodiles to better tolerate life in saline water and migrating through it. Because alligators and caimans have lost this ability, they are largely restricted to freshwater habitats, although larger alligators do sometimes live in tidal mangroves and in very rare cases in coastal areas.
Integumentary sense organs: Both crocodiles and alligators have small, pit-like sensory organs called integumentary sense organs (ISOs) or dermal pressure receptors (DPRs) surrounding their upper and lower jaws. These organs allow crocodilians to detect minor pressure changes in surrounding water, and assist them in locating and capturing prey. In crocodiles, however, such organs extend over nearly the entire body. Crocodile ISOs may also assist in detection of local salinity, or serve other chemosensory functions.
Less consistent differences: Crocodiles are generally thought of as more aggressive than alligators. Only six of the 23 crocodilian species are considered dangerous to adult human beings, most notably the Nile crocodile and saltwater crocodile. Each year, hundreds of deadly attacks are attributed to the Nile crocodile in sub-Saharan Africa. The American crocodile is considered to be less aggressive. Only a few (unverified) cases of American crocodiles fatally attacking humans have been reported.
Image gallery of extant species
| Biology and health sciences | Reptiles | null |
267376 | https://en.wikipedia.org/wiki/Mauna%20Kea%20Observatories | Mauna Kea Observatories | The Mauna Kea Observatories (MKO) are a group of independent astronomical research facilities and large telescope observatories that are located at the summit of Mauna Kea on the Big Island of Hawaiʻi, United States. The facilities are located in a 525-acre (212 ha) special land use zone known as the "Astronomy Precinct", which is located within the 11,228-acre (4,544 ha) Mauna Kea Science Reserve. The Astronomy Precinct was established in 1967 and is located on land protected by the Historical Preservation Act for its significance to Hawaiian culture. The presence and continued construction of telescopes is highly controversial due to Mauna Kea's centrality in native Hawaiian religion and culture, as well as for a variety of environmental reasons.
The location is near ideal because of its dark skies from lack of light pollution, good astronomical seeing, low humidity, high elevation of , position above most of the water vapor in the atmosphere, clean air, good weather and low latitude location.
Origin and background
Significantly predating the current observatories there is evidence of active astronomy on Mauna Kea in the 1901 Land Office Map of the Island of Hawaii showing the "Hawaii Astronomy Station" near the Mauna Kea summit.
After studying photos for NASA's Apollo program that contained greater detail than any ground-based telescope, Gerard Kuiper began seeking an arid site for infrared studies. While he first began looking in Chile, he also made the decision to perform tests in the Hawaiian Islands. Tests on Maui's Haleakalā were promising, but the mountain was too low in the inversion layer and often covered by clouds. On the "Big Island" of Hawaiʻi, Mauna Kea is considered the second-highest island mountain in the world. While the summit is often covered with snow, the air is extremely dry. Kuiper began looking into the possibility of an observatory on Mauna Kea. After testing, he discovered the low humidity was perfect for infrared signals. He persuaded Hawaiʻi Governor John A. Burns to bulldoze a dirt road to the summit where he built a small telescope on Puu Poliahu, a cinder cone peak. The peak was the second highest on the mountain with the highest peak being holy ground, so Kuiper avoided it. Next, Kuiper tried enlisting NASA to fund a larger facility with a large telescope, housing and other needed structures. NASA, in turn, decided to make the project open to competition. Professor of physics, John Jefferies of the University of Hawaii placed a bid on behalf of the university. Jefferies had gained his reputation through observations at Sacramento Peak Observatory. The proposal was for a two-meter telescope to serve both the needs of NASA and the university. While large telescopes are not ordinarily awarded to universities without well-established astronomers, Jefferies and UH were awarded the NASA contract, infuriating Kuiper, who felt that "his mountain" had been "stolen" from him. Kuiper would abandon his site (the very first telescope on Mauna Kea) over the competition and begin work in Arizona on a different NASA project. After considerable testing by Jefferies' team, the best locations were determined to be near the summit at the top of the cinder cones. Testing also determined Mauna Kea to be superb for nighttime viewing due to many factors, including the thin air, constant trade winds and being surrounded by sea. Jefferies would build a 2.24 meter telescope with the State of Hawaiʻi agreeing to build a reliable, all weather roadway to the summit. Building began in 1967 and first light was seen in 1970.
Other groups began requesting subleases on the newly accessible mountaintop. By 1970, two telescopes had been constructed by the United States Air Force and Lowell Observatory. In 1973, Canada and France agreed to build the 3.6 m CFHT on Mauna Kea. However, local organizations started to raise concerns about the environmental impact of the observatory. This led the Department of Land and Natural Resources to prepare an initial management plan, drafted in 1977 and supplemented in 1980. In January 1982, the University of Hawaiʻi Board of Regents approved a plan to support the continued development of scientific facilities at the site. In 1998, were transferred from the observatory lease to supplement the Mauna Kea Ice Age Reserve. The 1982 plan was replaced in 2000 by an extension designed to serve until 2020: it instituted an Office of Mauna Kea Management, designated for astronomy, and shifted the remaining to "natural and cultural preservation". This plan was further revised to address concern expressed in the Hawaiian community that a lack of respect was being shown toward the cultural value the mountain embodied to the region's indigenous people.
, the Mauna Kea Science Reserve has 13 observation facilities, each funded by as many as 11 countries. It is one of the world's premier observatories for optical, infrared, and submillimeter astronomy, and in 2009 was the largest measured by light gathering power. There are nine telescopes working in the visible and infrared spectrum, three in the submillimeter spectrum, and one in the radio spectrum, with mirrors or dishes ranging from . In comparison, the Hubble Space Telescope has a mirror, similar in size to the UH88, now the second smallest telescope on the mountain.
Controversies
Planned new telescopes, including the Thirty Meter Telescope, have attracted controversy due to their potential cultural and ecological impact. The multi-telescope "outrigger" extension to the Keck telescopes, which required new sites, was eventually canceled. Three or four of the mountain's 13 existing telescopes must be dismantled over the next decade with the TMT proposal to be the last area on Mauna Kea on which any telescope would ever be built.
With all the controversy, the building of telescopes has led to the creation of the Hawaii Night Sky Protection Act. As artificial light forms a light cloud above the land, the excess light disrupts the clear pictures taken by the telescopes. On July 1, 2013, the Hawaii Night Sky Protection Act was initiated, affecting both the Big Island and Maui. A large difference between the Hawaiian islands and the mainland United States can be observed: street lighting. Almost all streets are dark as the lamps have either not been built, have been removed, or have been dimmed.
Scientific discoveries
The Mauna Kea Observatories involves thirteen large telescopes. In November 2020, in collaboration with Europe's low-frequency ARray radio telescope, Mauna Kea's Gemini observatory and the NASA Infrared Telescope Facility (NITF), discovered the first Super-Planet.
In October, 2011, the Nobel Prize in Physics was awarded to Saul Perlmutter, Brian P. Schmidt, and Adam G. Riess; the prize recognized their findings, based on research on supernovae at the observatories, that dark energy is a force causing the universe to expand at an accelerating rate.
Management
The Reserve was established in 1968, and is leased by the State of Hawaiʻi's Department of Land and Natural Resources (DLNR). The University of Hawaiʻi manages the site and leases land to several multi-national facilities, which have invested more than $2 billion in science and technology. The lease expires in 2033 and after that 40 of 45 square kilometers (25 of 28 square miles) revert to the state of Hawaii.
Currently, pressure has been placed on the University based management. House Bill 2024 pushes for new stewardship on the summit. The area will be integrated in terms of management. Some land will be jurisdicted by the local community. The Bill pushes for future laws and regulations in terms of new building. As always, the cultural and environmental sites will be heavily recognized as a factor for consideration.
Location
The altitude and isolation in the middle of the Pacific Ocean makes Mauna Kea one of the best locations on Earth for ground-based astronomy. It is an ideal location for submillimeter, infrared and optical observations. The seeing statistics show that Mauna Kea is the best site in terms of optical and infrared image quality; for example, the CFHT site has a median seeing of 0.43 arcseconds.
Accommodations for research astronomers are located at the Onizuka Center for International Astronomy (often called Hale Pōhaku), by unpaved steep road from the summit at above sea level.
An adjacent visitor information station is located at . The summit of Mauna Kea is so high that tourists are advised to stop at the visitor station for at least 30 minutes to acclimate to atmospheric conditions before continuing to the summit, and scientists often stay at Hale Pōhaku for eight hours or more before spending a full night at observatories on the summit, with some telescopes requiring observers to spend one full night at Hale Pōhaku before working at the summit.
Telescopes
Telescopes found at the summit of Mauna Kea are funded by government agencies of various nations. The University of Hawaiʻi directly administers two telescopes. In total, there are twelve facilities housing thirteen telescopes at or around the summit of Mauna Kea.
Caltech Submillimeter Observatory (CSO): Caltech — closed 2015 — removed to be moved to Chile as the Leighton Chajnantor Telescope
Canada–France–Hawaiʻi Telescope (CFHT): Canada, France, University of Hawaiʻi
Gemini North Telescope: United States, United Kingdom, Canada, Chile, Australia, Argentina, Brazil
NASA Infrared Telescope Facility (IRTF): NASA
James Clerk Maxwell Telescope (JCMT): China, Japan, South Korea, Taiwan, United Kingdom, Canada
Subaru Telescope: National Astronomical Observatory of Japan
Sub-Millimeter Array (SMA): Taiwan, United States
Thirty Meter Telescope, proposed for Mauna Kea, alternatively the Canary Islands
United Kingdom Infrared Telescope (UKIRT): Lockheed Martin Advanced Technology Center, University of Hawaiʻi, University of Arizona
University of Hawaiʻi telescope (UH88): University of Hawaiʻi
University of Hawaiʻi telescope (Hoku Kea): University of Hawaii at Hilo
One receiver of the Very Long Baseline Array (VLBA): United States
W. M. Keck Observatory: California Association for Research in Astronomy
CSO, UKIRT, and Hoku Kea are scheduled for decommissioning as part of the Mauna Kea Comprehensive Management Plan.
Opposition and protests
In Honolulu, the governor and legislature, enthusiastic about the development, set aside an even larger area for the observatory after the initial project, causing opposition on the Big Island, in the city of Hilo. Native Hawaiians (kānaka ʻōiwi) believed the entire site was sacred and that developing the mountain, even for science, would spoil the area. Environmentalists were concerned about rare native bird populations and other citizens of Hilo were concerned about the sight of the domes from the city. Using town hall meetings, Jefferies was able to overcome opposition by weighing the economic advantage and prestige the island would receive. There has been substantial opposition to the Mauna Kea observatories that continues to grow. Over the years, the opposition to the observatories may have become the most visible example of the conflict science has encountered over access and use of environmental and culturally significant sites. Opposition to development grew shortly after expansion of the observatories commenced. Once access was opened up by the roadway to the summit, skiers began using it for recreation and objected when the road was closed as a precaution against vandalism when the telescopes were being built. Hunters voiced concerns, as did the Hawaiian Audubon Society who were supported by Governor George Ariyoshi.
The Audubon Society objected to further development on Mauna Kea over concerns to habitat of the endangered Palila, a species endemic to only specific parts of this mountain. The bird is the last of the finch billed honeycreepers existing on the island. Over 50% of native bird species had been killed off due to loss of habitat from early western settlers or the introduction of non-native species competing for resources. Hunters and sportsmen were concerned that the hunting of feral animals would be affected by the telescope operations. A "Save Mauna Kea" movement was inspired by the proliferation of telescopes, with opposition believing development of the mountain to be sacrilegious. Native Hawaiian non-profit groups, such as Kahea, whose goals are the protection of cultural heritage and the environment, oppose development on Mauna Kea as a sacred space to the Hawaiian religion. The land is protected by the United States Historical Preservation Act due to its significance to Hawaiian culture, but still allowed development.
2006 Kiholo Bay earthquake
A number of the telescopes sustained minor damage during the October 15, 2006 Kiholo Bay earthquake and aftershocks. JCMT was performing an inclinometry run and recorded the earthquake on its tilt sensors. Both CFHT and W. M. Keck Observatory were operational and back online by October 19.
Gallery
| Technology | Ground-based observatories | null |
267413 | https://en.wikipedia.org/wiki/P%C3%A8re%20David%27s%20deer | Père David's deer | The Père David's deer (Elaphurus davidianus), also known as the milu () or elaphure, is a species of deer native to the subtropical river valleys of China. It grazes mainly on grass and aquatic plants. It is the only extant member of the genus Elaphurus. Some experts suggest demoting Elaphurus to a subgenus of Cervus. Based on genetic comparisons, Père David's deer is closely related to Eld's deer.
Père David's deer were hunted almost to extinction in their native China by the late 19th century, but a number were taken to zoos in France and Germany and the deer was bred successfully in captivity. In the early 20th century, the British nobleman and politician Herbrand Russell, 11th Duke of Bedford, acquired a few Père David's deer from the Berlin Zoo and built up a large herd on his estate at Woburn Abbey. In the 1980s, the duke's great-grandson Robin Russell, 14th Duke of Bedford, donated several dozen deer to the Chinese government for reintroducing the species to the wild. As of 2020, the wild population in China was an estimated 2,825 individuals, with a further 7,380 in various nature reserves in China. All Père David's deer alive today descend from Herbrand Russell's original herd.
Demography
The Père David's deer is endemic to the Chinese region. According to fossil records, the species first appeared during the Pleistocene period, when it could be found across Manchuria. This demography changed during the Holocene period; during this time, the species could only be found in the swamp lands and wetlands of southern China. Due to hunting and land reclamation, the population of the Père David's deer became even smaller. By 1939, the last of the wild species were shot and killed.
Naming and etymology
This species of deer was first made known to Western science in 1866 by Armand David (Père David), a French missionary working in China. He obtained some hides and the carcasses of an adult male, an adult female and a young male, and sent them to Paris, where the species was named "Père David's deer" by Alphonse Milne-Edwards, a French biologist.
The species is sometimes known by its informal name sibuxiang (; Japanese: shifuzō), literally meaning "four not alike", which could mean "the four unlikes" or "like none of the four"; it is variously said that the four are cow, deer, donkey, horse (or) camel, and that the expression means in detail:
"the hooves of a cow but not a cow, the neck of a camel but not a camel, antlers of a deer but not a deer, the tail of a donkey but not a donkey."
"the nose of a cow but not a cow, the antlers of a deer but not a deer, the body of a donkey but not a donkey, tail of a horse but not a horse"
"the tail of a donkey, the head of a horse, the hoofs of a cow, the antlers of a deer"
"the neck of a camel, the hoofs of a cow, the tail of a donkey, the antlers of a deer"
"the antlers of a deer, the head of a horse and the body of a cow"
By this name, this undomesticated animal entered Chinese mythology as the mount of Jiang Ziya in Fengshen Bang (translated as Investiture of the Gods), a Chinese classical work of fiction written during the Ming dynasty.
Characteristics
The adult Père David's deer reaches a head-and-body length of up to and stands about tall at the shoulder. The tail is relatively long for a deer, measuring when straightened. Weight is between . The head is long and slender with large eyes, very large preorbital glands, a naked nose pad and small, pointed ears.
The branched antlers are unique in that the long tines point backward, while the main beam extends almost directly upward. There may be two pairs per year. The summer antlers are the larger set, and are dropped in November, after the summer rut. The second set—if they appear—are fully grown by January, and fall off a few weeks later.
The coat is reddish tan in the summer, changing to a dull gray in the winter. Long wavy guard hairs are present on and coat throughout the year, with the coat becoming woolier in winter. There is a mane on the neck and throat and a black dorsal stripe running along the cervicothoracic spine. The tail is about in length, with a dark tuft at the end. The hooves are large and spreading, and make clicking sounds (as in the reindeer) when the animal is moving.
A semiaquatic animal, Père David's deer swims well, spending long periods standing in water up to its shoulders. Although predominantly a grazer, the deer supplements its grass diet with aquatic plants in the summer.
Behavior
Père David's deer has similar reproductive physiological mechanisms to other deer species living in temperate latitudes. These mechanisms aid in the adaptation to a high-latitude environment. The reproductive behavior in stags differs from hinds.
In stags rutting behavior includes urine sniffing, anogenital sniffing, wallowing, and antler adorning. Communication behavior includes the spraying of urine and preorbital gland marking. The stags generally begin to rut before any signs of female sexual behavior. The reason to this is to establish a social rank among the stags. For Père David's deer the stag initiates the breeding season with the rutting.
In hinds, estrous behavior includes urinating frequently, receptivity, and allowing the stags to mount. Parental behavior involves sniffing the calf, calling, and rubbing the face of the calf. The hind has a typical behavior of being solitary.
Birth and lifespan
The gestation period is about nine months, roughly around 280 days, after which a single offspring is usually born; twins are born rarely. The gestation period is significantly longer than any other deer besides the roe deer. Père David's deer are considered seasonal breeders because three out of four calves are born during April in captive European populations. The breeding season is 160 days with the mating season usually being in June and July. Calf weight, calculated using Robbins and Robbins equation, tends to be between 11.3 and 13.2 kg. The juveniles (referred to as either fawns or calves) have a spotted coat, as is commonly seen in most species of deer. They are known to develop very quickly after birth. They reach sexual maturity at about 14 months. The average lifespan of a Père David deer is up to 18 years.
In captivity
An experiment was conducted to show how captivity would affect the deer's behavior and survival rate. Two areas were created; one with a large area containing a few Père David's deer present and one with a small area containing a high concentration of Père David's deer present. It was found when in captivity, it is best to keep the deer in large open areas that allow for adequate space with a reasonable number of individual deer living in the area; if put in a small area with a multitude of individuals present, stress amongst the deer will build up. Studies have shown through the high concentration of fecal matter in an area of captivity with limited space and a large number of Père David deer that they have shown different behavioral patterns to their wild counterparts. These deer would spend less time resting, and would stand longer due to human presence as well as display social aggression and competition over food sources. It is therefore unwise to keep them in small, densely populated areas to prevent accumulation of stress and aggression against each other.
The Père David's deer has been in isolation from the wild for more than 1200 years, causing humans to be the primary perceived threat due to a long-term lack of exposure to other natural predators. When encountering humans, the deer's response varies according to sex. Female-only groups display a lower degree of caution towards humans than in male-only groups. It was hypothesized that the presence of a single type of threat may be sufficient to maintain anti-predator responses in the face of relaxed predation pressure. During rutting season, the does will display increased vigilance and heightened threat perception, likely in response to the mating activities of the males.
Predation
Historically, the main predators of the Père David's deer are believed to have been tigers and leopards. Although they no longer encounter these predators while living in captivity, while experimenting with exposures to images and stimuli relating to these big cats, the deer seemed to instinctively react with a cautious predator–prey response typical of wild deer.
Population
In Neolithic times, the milus range extended across much of China proper. Archaeologists have found milu antlers at settlements from the Liao River in the north to Jiangsu and Zhejiang Province and across the Yellow and Yangtze River Basins in Shaanxi and Hunan Province. According to official data, the total population of the species has exceeded 8,000 in China. Milu is also under first-class state protection in China. In 1985, China Biodiversity Conservation and Green Development Foundation (CBCGDF) was established to help receive 22 Père David's deer from the Marquess of Tavistock of Woburn Abbey, England as a gift to return to their ancestral soil.
Extirpation in China
In the late 19th century, the world's only herd belonged to Tongzhi, the Emperor of China. The herd was maintained in the Nanyuan Royal Hunting Garden in Nan Haizi, near Peking. In 1895, one of the walls of the hunting garden was destroyed by a heavy flood of the Yongding River, and most of the deer escaped and some of them were killed and eaten by starving peasants. Fewer than thirty Père David's deer remained in the garden. Then in 1900, during the Boxer Rebellion, the garden was occupied by troops from the German Empire and all the remaining deer were shot and eaten by the soldiers, leaving the Père David's deer extirpated in its native China. A few of the deer had been legally obtained by the French and British Missions in Beijing and transported to various European zoos for exhibition and breeding. After the extirpation of the Chinese population in 1900, the English nobleman Herbrand Russell, 11th Duke of Bedford, was instrumental in saving the species. He acquired the few remaining deer from European zoos and formed a breeding herd in the deer park at his home at Woburn Abbey in Bedfordshire. Threatened again by both World Wars, the species survived largely due to the efforts of Bedford and his son Hastings, later 12th Duke of Bedford. The current world population, now found in zoos around the world, stems from the Woburn Abbey herd. Only three founders (1 male, two females) from the Berlin Zoo contributed to the current population of more than 5,000.
Reintroduction
Reintroduction of Père David's deer to China began in 1985, with a herd of 20 deer (5 males and 15 females). This was followed in 1987 by a second herd, consisting of 18 deer (all females). Both herds had been drawn from the Woburn Abbey herd and were donated by Robin Russell, Marquess of Tavistock (d. 2003; the future 14th Duke of Bedford), the 11th duke's great-grandson. In 2005 the Beijing authorities erected a statue of the 14th duke at Nan Haizi to mark the 20th anniversary of the Milu reintroduction. The transportation was sponsored by the World Wildlife Fund. The relic site of the Nanyuan (or Nan Haizi) Royal Hunting Garden in the southern suburbs of Beijing was chosen as the site of re-introduction, creating the Beijing Milu Park. The population in China expanded to around 2,000 in 2005.
A second re-introduction into China was conducted in 1986 where 36 Père David's deer were chosen from five UK zoological gardens with the bulk of the deer coming from Whipsnade Wild Animal Park in Bedfordshire. These deer were introduced into Dafeng Milu National Nature Reserve, near the Yellow Sea coast in eastern China. In 2006 the population at this Nature Reserve had reached around 950 with an average annual population increase of 17%.
Of 1993, 30 deer taken from the herd at Beijing Milu Park were released into the Tian'ezhou Milu National Nature Reserve (a.k.a. Shishou Milu National Nature Reserve) in , Shishou. These were followed by another 34 deer taken from Beijing Park and released into the Tian'ezhou Reserve. In a 1998 flood, a number of deer escaped from the reserve and have since been living and multiplying in the wild. As of 2015, the number of deer living in the wild is 700. The average annual population growth rate for Père David's deer in Tianezhou Nature Reserve was 22.2%.
In 2002, 30 deer taken from the herd at Beijing Milu Park and 20 from Dafeng Nature Reserve were released into the Yuanyang Yellow River Nature Reserve.
When the species was assessed for the IUCN Red List (1996), it was classified as "critically endangered" in the wild, under criterion "D'''": "[wild] population estimated to number less than 50 mature individuals". In October 2008, they were officially listed as extinct in the wild. Upon the status being updated in March 2016, it was clarified on this species' IUCN Red List page that it would remain listed as "Extinct in the Wild" until the reintroduced population proved long-term viability. Today, there are 53 herds of Père David's deer in China. Nine of the herds consist of 25 or under deer, and the remaining herds have under 10 deer. Due to the small population size a lack of genetic diversity is expected, but in spite of the small population size, the animals do not appear to suffer genetic problems from the genetic bottleneck. The captive population in China has increased in recent years, and the possibility remains that free-ranging populations can be reintroduced in the near future.
When reintroduced into their habitat, the species could face many problems due to years in captivity. Relaxed selection and reproduction with no environmental pressure may have taken place for Père David's deer, due to captive breeding, which can result in the loss of adaptive anti-predator behavior. It is possible that when fully released in the environment from captivity, after generations of offspring, the species could be unable to retain parasite-defense behaviors like grooming. If the population is reintroduced into the environment with no protection against tick infestation, then they can face major problems if not adapted to that environment. Père David's deer may also become confused by other predators such as tigers, since they are no longer adapted to them. A study was done on members of the species in captivity using the sounds of wild tigers roaring and domestic dogs barking, and the deer did not respond to the barking of the dogs, but hearing the tiger roar caused the deer to spend more time being cautious after hearing the sound, thus showing that the deer still retained ancestral memories of their previous predator, the tiger.
Red deer hybrid
Though New Zealand lacked its own terrestrial mammals, European settlers had introduced numerous species of deer into the land for the use of farming and hunting. Some of these deer species had crossbred in the wild, creating hybrids, which in turn were then utilized in deer farms based on their apparent genetic improvements. Alongside this discovery, deer farms began the practice of inciting hybridization in order to encourage genetic advancement. This includes a hybrid between Père David's deer and Red deer. These F1 hybrids are unique for several reasons. To begin with, both male and female offspring of this hybrid remain fertile, a rare prospect especially for species that have such a genetic distance. Both species differ in seasonal behaviors, gestation length, behavioral traits, morphology, maturity size, and disease resistance. Père David's deer is also unique in that its antlers are unlike any other deer in the world. However, due to having different seasonal behavior each species would enter their mating season at different intervals, thus preventing natural mating from occurring. In response, artificial insemination was employed on Red deer hinds with the semen from Père David's deer. These F1 hybrids did not share similar mating seasons with Père David's deer and as such were able to successfully mate with other red deer naturally. Three F1 hybrid stags successfully mated naturally in a period from 1989 to 1991 with 144 hinds and semen had been used to artificially inseminate 114 other Red deer hinds producing over 300 backcross hybrids.
Legend and cultural significance
According to Chinese legend, when the tyrant King Zhou of Shang ruled the land more than 3,000 years ago, a horse, a donkey, an ox and a deer went into a cave in the forest to meditate and on the day the King executed his minister Bigan, the animals awoke from their meditation and turned into humans. They entered society, learned of the King's heinous acts and wanted to take recourse against the King, who was powerful. So they transformed themselves into one creature that combined the speed of the horse, the strength of the ox, the donkey's keen sense of direction and the nimble agility of the deer. This new animal then galloped to the Kunlun Mountains to seek the advice of the Primeval Lord of Heaven. The Lord was astonished at the sight of a creature that had antlers of a deer, hooves of an ox, face of a horse and tail of a donkey. "It's unlike any of four creatures!" he exclaimed. Upon learning of the animal's quest, Lord gave his blessing and dispatched the creature to his disciple the sage Jiang Ziya, who was battling the King. Jiang Ziya rode the creature to victory over the King and helped found the Zhou dynasty. After fulfilling its vow, the milu settled in the lower reaches of the Yangtze River. The animal became a symbol of good fortune and was sought by later emperors who believed eating the meat of the milu would lead to everlasting life.
| Biology and health sciences | Deer | Animals |
267694 | https://en.wikipedia.org/wiki/Nervous%20tissue | Nervous tissue | Nervous tissue, also called neural tissue, is the main tissue component of the nervous system. The nervous system regulates and controls body functions and activity. It consists of two parts: the central nervous system (CNS) comprising the brain and spinal cord, and the peripheral nervous system (PNS) comprising the branching peripheral nerves. It is composed of neurons, also known as nerve cells, which receive and transmit impulses to and from it , and neuroglia, also known as glial cells or glia, which assist the propagation of the nerve impulse as well as provide nutrients to the neurons.
Nervous tissue is made up of different types of neurons, all of which have an axon. An axon is the long stem-like part of the cell that sends action potentials to the next cell. Bundles of axons make up the nerves in the PNS and tracts in the CNS.
Functions of the nervous system are sensory input, integration, control of muscles and glands, homeostasis, and mental activity.
Structure
Nervous tissue is composed of neurons, also called nerve cells, and neuroglial cells. Four types of neuroglia found in the CNS are astrocytes, microglial cells, ependymal cells, and oligodendrocytes. Two types of neuroglia found in the PNS are satellite glial cells and Schwann cells. In the central nervous system (CNS), the tissue types found are grey matter and white matter. The tissue is categorized by its neuronal and neuroglial components.
Components
Neurons are cells with specialized features that allow them to receive and facilitate nerve impulses, or action potentials, across their membrane to the next neuron. They possess a large cell body (soma), with cell projections called dendrites and an axon. Dendrites are thin, branching projections that receive electrochemical signaling (neurotransmitters) to create a change in voltage in the cell. Axons are long projections that carry the action potential away from the cell body toward the next neuron. The bulb-like end of the axon, called the axon terminal, is separated from the dendrite of the following neuron by a small gap called a synaptic cleft. When the action potential travels to the axon terminal, neurotransmitters are released across the synapse and bind to the post-synaptic receptors, continuing the nerve impulse.
Neurons are classified both functionally and structurally.
Functional classification:
Sensory neurons (afferent): Relay sensory information in the form of an action potential (nerve impulse) from the PNS to the CNS
Motor neurons (efferent): Relay an action potential out of the CNS to the proper effector (muscles, glands)
Interneurons: Cells that form connections between neurons and whose processes are limited to a single local area in the brain or spinal cord
Structural classification:
Multipolar neurons: Have 3 or more processes coming off the soma (cell body). They are the major neuron type in the CNS and include interneurons and motor neurons.
Bipolar neurons: Sensory neurons that have two processes coming off the soma, one dendrite and one axon
Pseudounipolar neurons: Sensory neurons that have one process that splits into two branches, forming the axon and dendrite
Unipolar brush cells: Are excitatory glutamatergic interneurons that have a single short dendrite terminating in a brush-like tuft of dendrioles. These are found in the granular layer of the cerebellum.
Neuroglia encompasses the non-neural cells in nervous tissue that provide various crucial supportive functions for neurons. They are smaller than neurons, and vary in structure according to their function.
Neuroglial cells are classified as follows:
Microglial cells: Microglia are macrophage cells that make up the primary immune system for the CNS. They are the smallest neuroglial cell.
Astrocytes: Star-shaped macroglial cells with many processes found in the CNS. They are the most abundant cell type in the brain, and are intrinsic to a healthy CNS.
Oligodendrocytes: CNS cells with very few processes. They form myelin sheaths on the axons of a neuron, which are lipid-based insulation that increases the speed at which the action potential, can travel down the axon.
NG2 glia: CNS cells that are distinct from astrocytes, oligodendrocytes, and microglia. They serve as the developmental precursors of oligodendrocytes.
Schwann cells: The PNS equivalent of oligodendrocytes, they help maintain axons and form myelin sheaths in the PNS.
Satellite glial cell: Line the surface of neuron cell bodies in ganglia (groups of nerve body cells bundled or connected together in the PNS)
Enteric glia: Found in the enteric nervous system, within the gastrointestinal tract.
Classification of tissue
In the central nervous system:
Grey matter is composed of cell bodies, dendrites, unmyelinated axons, protoplasmic astrocytes (astrocyte subtype), satellite oligodendrocytes (non-myelinating oligodendrocyte subtype), microglia, and very few myelinated axons.
White matter is composed of myelinated axons, fibrous astrocytes, myelinating oligodendrocytes, and microglia.
In the peripheral nervous system:
Ganglion tissue is composed of cell bodies, dendrites, and satellite glial cells.
Nerves are composed of myelinated and unmyelinated axons, Schwann cells surrounded by connective tissue.
The three layers of connective tissue surrounding each nerve are:
Endoneurium. Each nerve axon, or fiber is surrounded by the endoneurium, which is also called the endoneurial tube, channel or sheath. This is a thin, delicate, protective layer of connective tissue.
Perineurium. Each nerve fascicle containing one or more axons, is enclosed by the perineurium, a connective tissue having a lamellar arrangement in seven or eight concentric layers. This plays a very important role in the protection and support of the nerve fibers and also serves to prevent the passage of large molecules from the epineurium into a fascicle.
Epineurium. The epineurium is the outermost layer of dense connective tissue enclosing the (peripheral) nerve.
Function
The function of nervous tissue is to form the communication network of the nervous system by conducting electric signals across tissue. In the CNS, grey matter, which contains the synapses, is important for information processing. White matter, containing myelinated axons, connects and facilitates nerve impulse between grey matter areas in the CNS.
In the PNS, the ganglion tissue, containing the cell bodies and dendrites, contain relay points for nerve tissue impulses. The nerve tissue, containing myelinated axons bundles, carry action potential nerve impulses.
Clinical significance
Tumours
Neoplasms (tumours) in nervous tissue include:
Gliomas (glial cell tumors)
Oligoastrocytoma, Choroid plexus papilloma, Ependymoma, Astrocytoma (Pilocytic astrocytoma, Glioblastoma multiforme), Dysembryoplastic neuroepithelial tumour, Oligodendroglioma, Medulloblastoma, Primitive neuroectodermal tumor
Neuroepitheliomatous tumors
Ganglioneuroma, Neuroblastoma, Atypical teratoid rhabdoid tumor, Retinoblastoma, Esthesioneuroblastoma
Nerve sheath tumors
Neurofibroma (Neurofibrosarcoma, Neurofibromatosis), Schwannoma, Neurinoma, Acoustic neuroma, Neuroma
| Biology and health sciences | Tissues | null |
267933 | https://en.wikipedia.org/wiki/Breed | Breed | A breed is a specific group of breedable domestic animals having homogeneous appearance (phenotype), homogeneous behavior, and/or other characteristics that distinguish it from other organisms of the same species. In literature, there exist several slightly deviating definitions. Breeds are formed through genetic isolation and either natural adaptation to the environment or selective breeding, or a combination of the two. Despite the centrality of the idea of "breeds" to animal husbandry and agriculture, no single, scientifically accepted definition of the term exists. A breed is therefore not an objective or biologically verifiable classification but is instead a term of art amongst groups of breeders who share a consensus around what qualities make some members of a given species members of a nameable subset.
Another point of view is that a breed is consistent enough in type to be logically grouped together and when mated within the group produce the same type. When bred together, individuals of the same breed pass on these predictable traits to their offspring, and this abilityknown as "breeding true"is a requirement for a breed. Plant breeds are more commonly known as cultivars. The offspring produced as a result of breeding animals of one breed with other animals of another breed are known as crossbreeds or mixed breeds. Crosses between animal or plant variants above the level of breed/cultivar (i.e. between species, subspecies, botanical variety, even different genera) are referred to as hybrids.
Breeding: selection by breeders
The breeder (or group of breeders) who initially establishes a breed does so by selecting individual animals from within a gene pool that they see as having the necessary qualities needed to enhance the breed model they are aiming for. These animals are referred to as foundation stock. Furthermore, the breeder mates the most desirable representatives of the breed from his or her point of view, aiming to pass such characteristics to their progeny. This process is known as selective breeding. A written description of desirable and undesirable breed representatives is referred to as a breed standard.
Breed characteristics
Breed specific characteristics, also known as breed traits, are inherited, and purebred animals pass such traits from generation to generation. Thus, all specimens of the same breed carry several genetic characteristics of the original foundation animal(s). In order to maintain the breed, a breeder would select those animals with the most desirable traits to achieve further maintenance and developing of such traits. At the same time, the breed would avoid animals carrying characteristics undesirable or not typical for the breed, including faults or genetic defects. The population within the same breed should consist of a sufficient number of animals to maintain the breed within the specified parameters without the necessity of forced inbreeding.
Domestic animal breeds commonly differ from country to country, and from nation to nation. Breeds originating in a certain country are known as "native breeds" of that country.
Lists of breeds
| Technology | Animal husbandry | null |
267968 | https://en.wikipedia.org/wiki/Tyranni | Tyranni | The Tyranni (suboscines) are a suborder of passerine birds that includes more than 1,000 species, a large majority of which are South American. It is named after the type genus Tyrannus. These have a different anatomy of the syrinx musculature than the oscines (songbirds of the larger suborder Passeri), hence the common name of suboscines.
The suboscines originated in South America about 50 million years ago and dispersed into the Old World likely via a trans-Atlantic route during the Oligocene. Their presence in the early Oligocene of Europe is well documented by several fossil specimens.
Systematics
The suborder Tyranni is divided into two infraorders: the Eurylaimides and the Tyrannides. The New Zealand wrens in the family Acanthisittidae are placed in a separate suborder Acanthisitti.
The phylogenetic relationships of the 16 families in the Tyranni suborder is shown below. The cladogram is based on a large molecular genetic study by Carl Oliveros and collaborators published in 2019: The families and the species numbers are from the list maintained by the International Ornithologists' Union (IOC).
The Eurylaimides contain the Old World suboscines – mainly distributed in tropical regions around the Indian Ocean – and a single American species, the sapayoa:
Philepittidae: asities
Eurylaimidae: typical broadbills
Calyptomenidae: African and green broadbills
Sapayoidae: broad-billed sapayoa
Pittidae: pittas
The Tyrannides contain all the suboscines from the Americas except the broad-billed sapayoa. The families listed here are those recognised by the International Ornithologists' Union.
Pipridae: manakins
Cotingidae: cotingas
Tityridae: tityras, sharpbill, becards (includes Oxyruncus and Onychorhynchus)
Tyrannidae: tyrant-flycatchers (includes Piprites, Platyrinchus, Tachuris and Rhynchocyclus)
Melanopareiidae: crescent chests
Conopophagidae: gnateaters and gnatpittas
Thamnophilidae: antbirds
Grallariidae: antpittas
Rhinocryptidae: tapaculos
Formicariidae: antthrushes
Furnariidae: ovenbirds and woodcreepers (includes Dendrocolaptidae)
This group has been separated into three parvorders by Sibley & Ahlquist. However, DNA:DNA hybridization did not reliably resolve the suboscine phylogeny. It was eventually determined that there was a simple dichotomy between the antbirds and allies (tracheophones), and the tyrant-flycatchers and allies. Given that the "parvorder" arrangement originally advanced is obsolete (see e.g. Irestedt et al. 2002 for tracheophone phylogeny) — more so if the Eurylaimides are elevated to a distinct suborder — it is better to rank the clades as superfamilies or, if the broadbill group is considered a separate suborder, as infraorders. In the former case, the name Furnarioidea would be available for the tracheophones, whereas "Tyrannoidea", the "bronchophone" equivalent, has not yet been formally defined.
In the latter case, the tracheophones would be classified as "Furnariides",
while the Tyrannides would be restricted to the tyrant-flycatchers and other "bronchophone" families.
The tracheophones contain the Furnariidae, Thamnophilidae, Formicariidae (probably including most tapaculos), and Conopophagidae. The tyrant-flycatcher clade includes the namesake family, the Tityridae, the Cotingidae, and the Pipridae.
| Biology and health sciences | Tyranni | Animals |
268115 | https://en.wikipedia.org/wiki/Ocelot | Ocelot | The ocelot (Leopardus pardalis) is a medium-sized spotted wild cat that reaches at the shoulders and weighs between on average. It is native to the southwestern United States, Mexico, Central and South America, and the Caribbean islands of Trinidad and Margarita. Carl Linnaeus scientifically described it in 1758. Two subspecies are recognized.
The ocelot is efficient at climbing, leaping and swimming. It prefers areas close to water sources with dense vegetation cover and high prey availability. It preys on small terrestrial mammals, such as armadillos, opossums, and lagomorphs. It is typically active during twilight and at night and tends to be solitary and territorial. Both sexes become sexually mature at around two years of age and can breed throughout the year; peak mating season varies geographically. After a gestation period of two to three months, the female gives birth to a litter of one to three kittens. They stay with their mother for up to two years, after which they leave to establish their own home ranges.
The ocelot is listed as Least Concern on the IUCN Red List and is threatened by habitat destruction, hunting, and traffic accidents. While its range is very large, various populations are decreasing in many parts of its range. The association of the ocelot with humans dates back to the Aztec and Incan civilizations; it has occasionally been kept as a pet.
Etymology
The name "ocelot" comes from the Nahuatl word (), which generally refers to the jaguar, rather than the ocelot. Another possible origin for the name is the Latin ("having little eyes" or "marked with eye-like spots"), in reference to the cat's spotted coat.
Other vernacular names for the ocelot include (Venezuela), (Argentina), (Panama), (Suriname), , (Brazil), (Costa Rica, Panama and Venezuela), , , , , (Belize), (Bolivia) and (Colombia, Ecuador, Guatemala, and Peru).
Taxonomy
Felis pardalis was the scientific name proposed for the ocelot by Carl Linnaeus in 1758. The genus Leopardus was proposed by John Edward Gray in 1842 for several spotted cat skins in the collection of the Natural History Museum, London.
Several ocelot specimens were described in the nineteenth and twentieth centuries, including:
Felis mitis by Frédéric Cuvier in 1824 was a specimen from Rio de Janeiro, Brazil.
F. chibi-gouazou by Edward Griffith in 1827 was based on earlier descriptions and illustrations.
Leopardus griseus by John Edward Gray in 1842 was a spotted cat skin from Central America.
F. pseudopardalis by Pierre Boitard in 1845 was an ocelot kept in the Jardin des plantes.
F. melanura by Robert Ball in 1844 was a specimen from British Guiana.
F. albescens by Jacques Pucheran in 1855 was a specimen from Brownsville, Texas.
F. aequatorialis by Edgar Alexander Mearns in 1903 was a skin of an adult female ocelot from Talamanca canton in Costa Rica.
F. maripensis and F. sanctaemartae by Joel Asaph Allen in 1904 were skins of two adult female ocelots from Maripa, Venezuela and Santa Marta district in Colombia, respectively.
F. pardalis pusaea by Oldfield Thomas in 1914 was an ocelot skin and skull from Guayas Province in coastal Ecuador.
F. pardalis nelsoni and F. p. sonoriensis by Edward Alphonso Goldman in 1925 as subspecies of F. pardalis, based on specimens from Manzanillo and the Mayo River region respectively in Mexico.
L. pardalis steinbachi by Reginald Innes Pocock in 1941 was a specimen from Buena Vista, Ichilo in Bolivia.
Subspecies
In 1919, Allen reviewed the specimens described until 1914, placed them into the genus Leopardus and recognized nine subspecies as valid taxa based on the colors and spot patterns of skins. In 1941, Pocock reviewed dozens of ocelot skins in the collection of the Natural History Museum and regrouped them to nine different subspecies, also based on their colors and spots. Later authors recognized 10 subspecies as valid.
In 1998, results of a mtDNA control region analysis of ocelot samples indicated that four major ocelot groups exist, one each in Central America, northwestern South America, northeastern South America and southern South America south of the Amazon River. A 2010 study of morphological features noted significant differences in the size and color of the Central and South American populations, suggesting they could be separate species. In 2013, a study of craniometric variation and microsatellite diversity in ocelots throughout the range recognized three subspecies: L. p. albescens from the Texas–Mexico border, L. p. pardis from Central America and L. p. pseudopardalis from South America, though L. p. mitis may comprise the ocelot population in the southern part of South America.
In 2017, the Cat Classification Task Force of the IUCN Cat Specialist Group noted that up to four subspecies can be identified, but recognized only two as valid taxa. These two taxa differ in morphological features and are geographically separated by the Andes:
L. p. pardalis has a greyish fur. Its range extends from Texas and Arizona to Costa Rica.
L. p. mitis has a more yellowish fur and is larger than pardalis. It occurs in South America as far south as northern Argentina.
Phylogeny
Results of a phylogenetic study indicate that the Leopardus lineage genetically diverged from the Felidae around 8 million years ago (mya). The ocelot is estimated to have diverged from the margay (Leopardus wieldii) between 2.41 and 1.01 mya. The relationships of the ocelot within the Felidae is considered as follows:
Characteristics
The ocelot's fur is extensively marked with solid black markings on a creamy, tawny, yellowish, reddish gray or gray background color. The spots on the head and limbs are small, but markings on the back, cheeks, and flanks are open or closed bands and stripes. A few dark stripes run straight from the back of the neck up to the tip of the tail. Its neck and undersides are white, and the insides of the legs are marked with a few horizontal streaks. Its round ears are marked with a bright white spot. Its fur is short, about long on the belly, but with about long guard hairs on the back. The body has a notably strong odor. Each ocelot has a unique color pattern, which can be used to identify individuals. Its eyes are brown, but reflect in a golden hue when illuminated. It has 28 to 30 teeth, with the dental formula . It has a bite force quotient at the canine tip of 113.8. Only one ocelot is known to possess albinism, and the appearance of such a trait in ocelots is likely an indication of shrinking populations due to deforestation.
With a head-and-body length ranging from and a long tail, the ocelot is the largest member of the genus Leopardus. It typically reaches at the shoulder. The weight of females ranges between and of males between . Its footprint measures nearly .
The ocelot can be confused with the margay (Leopardus wiedii) and the oncilla (L. tigrinus), though the ocelot is noticeably larger and heavier with a shorter tail. Though all three have rosettes on their coats, the ocelot typically has a more blotched pattern; the oncilla has dark spots on its underbelly unlike the other two. Other differences lie in the facial markings, appearance of the tail and fur characteristics. The ocelot is similar in size to a bobcat (Lynx rufus), though larger individuals have occasionally been recorded. The jaguar is notably larger and heavier, and has rosettes instead of spots and stripes.
Distribution and habitat
The ocelot ranges from the southwestern United States to northern Argentina, up to an elevation of . In the United States, it occurs in Texas and Arizona, and is extirpated from Louisiana and Arkansas. Ocelots fossils were found in Florida.
It inhabits tropical forests, thorn forests, mangrove swamps and savannas. In the Amazon rainforest, it prefers habitats with availability of prey and water, and tends to avoid other predators. It favors areas with dense forest cover and water sources, far from roads and human settlement, avoiding steep slopes and highly elevated areas. In areas where ocelots coexist with larger predators such as cougars and humans, they tune their active hours to avoid them, and seek dense cover to avoid competitors. It can adapt well to its surroundings; as such, factors other than the aforementioned are not significant in its choice of habitat.
It shares a large part of its range with the jaguar, jaguarundi, margay, oncilla and cougar.
Ecology and behavior
The ocelot is usually solitary and active mainly during twilight and at night. Radio collared individuals in the Cocha Cashu Biological Station in Peru rested during the day and became active earliest in the late afternoon; they moved between 3.2 and 17 hours until dawn and then returned to their dens.
During the daytime, it rests on trees, in dens below large trees or other cool, sheltered sites on the ground. It is agile in climbing and leaping, and escapes predators by jumping on trees. It is also an efficient swimmer. It scent-marks its territory by spraying urine. The territories of males are large, while those of females cover . Territories of females rarely overlap, whereas the territory of a male includes those of two to three females. Social interaction between sexes is minimal, though a few adults have been observed together even in non-mating periods, and some juveniles interact with their parents. Data from camera trapping studies confirm that several ocelot individuals deposit scat in one or several communal sites, called latrines. Ocelots can be aggressive in defending their territory, fighting even to death.
The population density of ocelots has been observed to be high in areas with high rainfall, but tends to decrease with increasing latitude; highest densities have been recorded in the tropics. In 2014, the ocelot population density in Barro Colorado Island was estimated to be , greater than recorded in northwestern Amazon in Peru in 2010, which was the densest ocelot population recorded thus far.
Potential predators of the ocelot in Texas include the cougar, coyote and American alligator, while ocelot kittens are vulnerable to raptors, such as the great horned owl, as well as feral dogs, feral pigs and snakes. Studies have found that adult ocelots are vulnerable to predation by both cougars and jaguars, with decreasing water sources in Guatemala causing predatory encounters with the latter.
Hunting and diet
Ocelots have been observed to follow scent trails in search for prey, walking at a speed of about . Alternatively, an ocelot may wait for prey for 30 to 60 minutes at a certain site and move to another walking at if unsuccessful. An ocelot typically prefers hunting in areas with vegetation cover, avoiding open areas, especially on moonlit nights, so as not to be seen by the prey. As a carnivore, it preys on small terrestrial mammals such as rodents, lagomorphs, armadillos, opossums, also fish, crustaceans, insects, reptiles and birds. It usually feeds on the kill immediately, but removes bird feathers before. It typically preys on animals that weigh less than , but rarely targets large ungulates such as deer, sheep and peccaries, as well as anteaters, New World monkeys and iguanas. It requires of food every day to satisfy its energy requirements.
Primates prevail in the diet of ocelots in southeastern Brazil and iguanas in a tropical deciduous forest in Mexico. The composition of the diet varies by season; in Venezuela, ocelots were found to prefer iguanas and rodents in the dry season and then switch to land crabs in the wet season. In southeastern Brazil, ocelots have a similar prey preference as margays and oncillas. The oncillas focus on tree-living marsupials and birds while the margays are not as selective.
Reproduction and life cycle
Both male and female ocelots produce a long-range "yowl" in the mating season and a short-range "meow". Ocelots can mate any time during the year. The peak mating season varies geographically; in Argentina and Paraguay, peaks have been observed in autumn, in Mexico and Texas in autumn and winter. Estrus lasts four to five days and recurs every 25 days in a non-pregnant female. A study in southern Brazil showed that sperm production in ocelots, margays and oncillas peaks in summer. When mating, captive ocelots spend more time together, scent-mark extensively and eat less.
The female gives birth to a litter of one to three kittens after a gestation period of two to three months. Dens are usually located in dense vegetation. A newborn kitten weighs . The kitten is born with spots and stripes, though on a gray background; the color changes to golden as the ocelot grows older. A study in southern Texas revealed that a mother keeps a litter in a den for 13 to 64 days and shifts the young to two or three dens. The kitten's eyes open 15 to 18 days after birth. Kittens begin to leave the den at the age of three months. They remain with their mother for up to two years and then start dispersing and establishing their own territory. In comparison to other felids, ocelots have a relatively longer duration between births and a narrow litter size. Captive ocelots live for up to 20 years.
Threats
Throughout its range, the ocelot is threatened by loss and fragmentation of habitat. In Texas, the fertile land that supports dense cover and constitutes the optimum habitat for the ocelot is being lost to agriculture. The habitat is often fragmented into small pockets that cannot support ocelots well, leading to deaths due to starvation. Traffic accidents have emerged as a major threat over the years, as ocelots try to expand beyond their natural habitat to new areas and get hit by vehicles. In the Atlantic Forest in northeastern Argentina, it is affected by logging and poaching of prey species.
The fur trade was a flourishing business in the 1960s and the 1970s that resulted in severe exploitation of felids such as the ocelot and the jaguar. In the 1960s, ocelot skins were among the most highly preferred in the US, reaching an all-time high of 140,000 skins traded in 1970. This was followed by prohibitions on commercial trade of spotted cat skins in several range states such as Brazil and the US, causing ocelot skins in trade to plummet. In 1986, the European Economic Community banned import of ocelot skins, and in 1989, the ocelot was included in Appendix I of the Convention on International Trade in Endangered Species of Wild Fauna and Flora. However, hunting of ocelots for skins has continued and is still a major threat to ocelot survival.
Another threat has been the international pet trade; this typically involves capturing ocelot kittens by killing their mothers; these cats are then sold to tourists. Though it is banned in several countries, pet trade survives; in some areas of Central and South America, ocelots are still sold in a few local markets.
Conservation
The ocelot is listed as Least Concern on the IUCN Red List because of its wide distribution in the Americas. Ocelot hunting is banned in Argentina, Brazil, Bolivia, Colombia, Costa Rica, French Guiana, Guatemala, Honduras, Mexico, Nicaragua, Panama, Paraguay, Suriname, Trinidad and Tobago, the United States, Uruguay, and Venezuela; hunting is regulated in Peru. As of 2013, the global population was estimated at more than 40,000 mature individuals. Ocelot populations were stable in some Amazon basin areas as of 2013. As of 2012, the ocelot population in Argentina's subtropical regions was estimated to consist of 1,500 to 8,000 mature individuals. It has been recorded in oil palm landscapes and big cattle ranches in the Colombian Llanos and inter-Andean valleys.
In Texas
In Texas and northeastern Mexico, ocelot populations have reduced drastically; as of 2014, the population in Texas was estimated to be 50–80 individuals. The reduced numbers have led to increased inbreeding and low genetic diversity. Despite this, the US Fish and Wildlife Service failed to acknowledge the ocelot population in Texas as a distinct population segment worthy of listing as endangered. The US Fish and Wildlife Service, the Texas Parks and Wildlife Department and The Nature Conservancy are among agencies actively involved in ocelot conservation efforts, such as the protection and regeneration of vegetation in the Rio Grande Valley. Much of the reintroduction effort is taking place on private lands. NatureServe considers the ocelot apparently secure globally, but critically imperiled in Texas and Arizona.
In captivity
The American Zoo and Aquarium Association established a Species Survival Plan for the ocelot populations in Brazil. In 2006, the captive population in North American zoos consisted of 16 ocelots representing six founders and their offspring. Some litters were produced using artificial insemination.
The Emperor Valley Zoo in Trinidad keeps foremost confiscated and trapped ocelots.
In culture
Ocelots have been associated with humans since the time of the Aztec and Incan civilizations, who depicted ocelots in their art and mythology. Representations of ocelots appear in every artistic medium, from Moche ceramics to murals, architectural details, and landscape features. Ocelot bones were made into thin, pointed instruments to pierce ears and limbs for ritual bloodletting. Several figurines depicting ocelots and similar felids are known. In her 1904 work A Penitential Rite of the Ancient Mexicans, archaeologist Zelia Nuttall described a statue depicting an ocelot or another felid excavated in Mexico City and its relation to the Aztec deity Tezcatlipoca. She argued that the sculpture depicted an ocelot, writing,
Moreover, she described a photograph of a seated person to corroborate her claim:
Like many other felids, occasionally ocelots are kept as pets. They might demand a lot of attention from their owners and have a tendency to chew on or suck on objects, such as fabric and the fingers of their owners; this can lead them to accidentally ingest objects such as tennis balls. Agile and playful, pet ocelots can be troublesome to keep due to their habit of leaping around and potentially damaging objects; ocelots may unintentionally injure their owners with bites. Nevertheless, carefully raised ocelots can be highly affectionate. Painter Salvador Dalí kept a pet ocelot named Babou that was seen with him at many places he visited, including a voyage aboard SS France. When one of the diners at a New York restaurant was alarmed by his ocelot, Dali told her that it was a common domestic cat that he had "painted over in an op art design". Opera singer Lily Pons and musician Gram Parsons are also known to have kept ocelots.
| Biology and health sciences | Felines | Animals |
268216 | https://en.wikipedia.org/wiki/Fuchsia | Fuchsia | Fuchsia ( ) is a genus of flowering plants that consists mostly of shrubs or small trees.
Almost 110 species of Fuchsia are recognized; the vast majority are native to South America, but a few occur north through Central America to Mexico, and also several from New Zealand to Tahiti. One species, F. magellanica, extends as far as the southern tip of South America, occurring on Tierra del Fuego in the cool temperate zone, but the majority are tropical or subtropical.
Taxonomy
The first to be scientifically described, Fuchsia triphylla, was discovered on the Caribbean island of Hispaniola (Haiti and the Dominican Republic) about 1696–1697 by the French Minim friar and botanist, Charles Plumier, during his third expedition to the Greater Antilles. He named the new genus after German botanist Leonhart Fuchs (1501–1566, ).
The fuchsias are most closely related to the northern hemisphere genus Circaea, the two lineages having diverged around 41 million years ago.
Description
Most fuchsias are shrubs from tall, but one New Zealand species, the kōtukutuku (F. excorticata), is unusual in the genus in being a tree, growing up to tall. Fuchsia leaves are opposite or in whorls of three to five, simple lanceolate, and usually have serrated margins (entire in some species), 1–25 cm long, and can be either deciduous or evergreen, depending on the species.
The flowers are very decorative; they have a pendulous teardrop shape and are displayed in profusion throughout the summer and autumn, and all year in tropical species. They have four long, slender sepals and four shorter, broader petals; in many species, the sepals are bright red and the petals purple (colours that attract the hummingbirds that pollinate them), but the colours can vary from white to dark red, purple-blue, and orange. A few have yellowish tones. The ovary is inferior.
The fruit is a small (5–25 mm) dark reddish green, deep red, or deep purple berry, containing numerous very small seeds. The fruit of the berry of F. splendens is reportedly among the best-tasting. Its flavor is reminiscent of citrus and black pepper, and it can be made into jam. The fruits of some other fuchsias are flavorless or leave a bad aftertaste.
Species
The majority of Fuchsia species are native to Central and South America. A small additional number are found on Hispaniola (two species), in New Zealand (three species) and on Tahiti (one species). Philip A. Munz in his A Revision of the Genus Fuchsia classified the genus into seven sections of 100 species. More recent scientific publications, especially those by the botanists Dennis E. Breedlove of the University of California and, currently, Paul E. Berry of the University of Michigan, recognize 108 species and 122 taxa, organized into 12 sections. In New Zealand and Tahiti, section Skinnera now consists of only three species as F. × colensoi has been determined to be a naturally occurring hybrid between F. excorticata and F. perscandens. Also, F. procumbens has been placed into its own section, Procumbentes. Two other new sections are Pachyrrhiza and Verrucosa, each with one species. The Plant List, a cooperative endeavor by several leading botanical institutions to maintain a working list of all plant species, lists most currently accepted Fuchsia species and synonyms.
The vast majority of garden hybrids have descended from a few parent species.
Section Ellobium
Mexico and Costa Rica. This section contains three species.
Fuchsia decidua
Fuchsia fulgens
Fuchsia splendens
Section Encliandra
Mexico to Panama. Flowers on the six species in this section have flat petals and short stamens and are reflexed into the tube. Fruits contain few seeds.
Fuchsia encliandra
Fuchsia encliandra subsp. encliandra
Fuchsia encliandra subsp. microphyloides
Fuchsia encliandra subsp. tetradactyla
Fuchsia microphylla
Fuchsia microphylla subsp. aprica
Fuchsia microphylla subsp. chiapensis
Fuchsia microphylla subsp. hemsleyana
Fuchsia microphylla subsp. hidalgensis
Fuchsia microphylla subsp. microphylla
Fuchsia microphylla subsp. quercertorum
Fuchsia obconica
Fuchsia parviflora
Fuchsia ravenii
Fuchsia thymifolia
Fuchsia thymifolia subsp. minimiflora
Fuchsia thymifolia subsp. thymiflora
Fuchsia × bacillaris
Section Fuchsia
Northern Argentina to Colombia and Venezuela, and Hispaniola. With 64 currently recognized species, Sect. Fuchsia (syn. Eufuchsia) is the largest section within the genus. The flowers are perfect, with convolute petals. The stamens are erect and may or may not be exserted from the corolla; the stamens opposite the petals are shorter. The fruit has many seeds.
Fuchsia abrupta
Fuchsia ampliata
Fuchsia andrei
Fuchsia aquaviridis
Fuchsia austromontana
Fuchsia ayavacensis
Fuchsia boliviana
Fuchsia campii
Fuchsia canescens
Fuchsia caucana
Fuchsia ceracea
Fuchsia cinerea
Fuchsia cochabambana
Fuchsia confertifolia
Fuchsia coriacifolia
Fuchsia corollata
Fuchsia corymbiflora
Fuchsia crassistipula
Fuchsia cuatrecasaii
Fuchsia decussata
Fuchsia denticulata
Fuchsia dependens
Fuchsia ferreyrae
Fuchsia fontinalis
Fuchsia furfuracea
Fuchsia gehrigeri
Fuchsia glaberrima
Fuchsia harlingii
Fuchsia hartwegii
Fuchsia hirtella
Fuchsia hypoleuca
Fuchsia lehmannii
Fuchsia llewelynii
Fuchsia loxensis
Fuchsia macrophylla
Fuchsia macropetala
Fuchsia macrostigma
Fuchsia magdalenae
Fuchsia mathewsii
Fuchsia nigricans
Fuchsia orientalis
Fuchsia ovalis
Fuchsia pallescens
Fuchsia petiolaris
Fuchsia pilosa
Fuchsia polyantha
Fuchsia pringsheimii
Fuchsia putumayensis
Fuchsia rivularis
Fuchsia rivularis subsp. pubescens
Fuchsia rivularis subsp. rivularis
Fuchsia sanctae-rosae
Fuchsia sanmartina
Fuchsia scabriuscula
Fuchsia scherffiana
Fuchsia sessifolia
Fuchsia simplicicaulis
Fuchsia steyermarkii
Fuchsia summa
Fuchsia sylvatica
Fuchsia tincta
Fuchsia triphylla
Fuchsia vargasiana
Fuchsia venusta
Fuchsia vulcanica
Fuchsia wurdackii
Section Hemsleyella
Venezuela to Bolivia. The fifteen species in this section are characterised by a nectary that is fused with the base of the flower tube and petals that are partly or completely lacking.
Fuchsia apetala
Fuchsia cestroides
Fuchsia chloroloba
Fuchsia garleppiana
Fuchsia huanucoensis
Fuchsia inflata
Fuchsia insignis
Fuchsia juntasensis
Fuchsia membranaceae
Fuchsia mezae
Fuchsia nana
Fuchsia pilaloensis
Fuchsia salicifolia
Fuchsia tilletiana
Fuchsia tunariensis
Section Jimenezia
Panama and Costa Rica.
Fuchsia jimenezii
Section Kierschlegeria
Coastal central Chile. This section is made up of a single species with pendulous axillary pedicels. The leaves are sparse. The sepals are reflexed and slightly shorter than the tube.
Fuchsia lycioides
Section Pachyrrhiza
Peru.
Fuchsia pachyrrhiza
Section Procumbentes
New Zealand.
Fuchsia procumbens
Section Quelusia
Southern Argentina and Chile, and Southeastern Brazil. The nine species in this section have the nectary fused to the base of the tube, or hypanthium. The hypanthium is cylindrical and is generally no longer than the sepals. The stamens are long and are exserted beyond the corolla.
Fuchsia alpestris
Fuchsia bracelinae
Fuchsia brevilobis
Fuchsia campos-portoi
Fuchsia coccinea
Fuchsia glazioviana
Fuchsia hatschbachii
Fuchsia magellanica
Fuchsia regia
Fuchsia regia subsp. regia
Fuchsia regia subsp. reitzii
Fuchsia regia subsp. serrae
Section Schufia
Mexico to Panama. These two species bear flowers in an erect, corymb-like panicle.
Fuchsia arborescens
Fuchsia paniculata
Fuchsia paniculata subsp. mixensis
Fuchsia paniculata subsp. paniculata
The name Schufia is a taxonomic anagram derived from Fuchsia.
Section Skinnera
New Zealand and Tahiti. The three living species have a floral tube with a swelling above the ovary. The sepals curve back on themselves and the petals are small or nearly absent. A new fossil species from the Early Miocene in New Zealand was described in October 2013.
†Fuchsia antiqua
Fuchsia cyrtandroides
Fuchsia excorticata
Fuchsia perscandens
Fuchsia × colensoi – a natural hybrid
Section Verrucosa
Venezuela and Colombia.
Fuchsia verrucosa
Cultivation
Fuchsias are popular garden shrubs, and once planted can live for years with a minimal amount of care. The British Fuchsia Society maintains a list of hardy fuchsias that have been proven to survive a number of winters throughout Britain and to be back in flower each year by July. Enthusiasts report that hundreds and even thousands of hybrids survive and prosper throughout Britain. In the United States, the Northwest Fuchsia Society maintains an extensive list of fuchsias that have proven hardy in members' gardens in the Pacific Northwest over at least three winters.
Fuchsias from sections Quelusia (F. magellanica, F. regia), Encliandra, Skinnera (F. excorticata, F. perscandens) and Procumbentes (F. procumbens) have especially proven to be hardy in widespread areas of Britain and Ireland, as well as in many other countries such as New Zealand (aside from its native species) or the Pacific Northwest region of the United States. A number of species will easily survive outdoors in agreeable mild temperate areas. Though some may not always flower in the average British summer, they will often perform well in other favorable climatic zones. Even in somewhat colder regions, a number of the hardier species will often survive as herbaceous perennials, dying back and reshooting from below ground in the spring.
Due to the favorably mild, temperate climate created by the North Atlantic Current fuchsias grow abundantly in the West Kerry and West Cork region of Ireland and in the Isles of Scilly, even colonising wild areas there. While F. magellanica is not widespread in Scotland it has been known to grow wild in sheltered areas, such as the banks of local streams in Fife. In the Pacific Northwest region of the United States, F. magellanica also easily survives regional winters.
Categories
Horticultural fuchsias may be categorised as upright and bushy, or trailing. Some can be trained as hedges, such as F. magellanica. Faster-growing varieties are easiest to train. Care should be taken to choose the hardier cultivars for permanent plantings in the garden as many popular upright Fuchsias such as 'Ernie', 'Jollies Nantes' and 'Maria Landy' are not reliably winter hardy, but rather extremely tender (hardiness zone 10).
Cultivars
In the UK, 60 cultivated varieties of fuchsia have gained the Royal Horticultural Society's Award of Garden Merit, including:-
'Alice Hoffman' (pink sepals, white petals – hardy)
'Display' (deep pink/rose pink single)
'Dollar Prinzessin' (cerise sepals, purple petals – hardy)
'Garden News' (light pink sepals, double magenta petals – hardy)
'Genii' (single, cerise/purple)
'Hawkshead' (white self)
'Lady Thumb' (compact, pink sepals, white petals)
'Mrs Popple' (vigorous, red sepals, purple petals – hardy)
'Riccartonii' (crimson sepals, purple petals)
'Snowcap' (scarlet sepals, white double petals)
'Swingtime' (double, scarlet sepals, white petals)
'Thalia' (tryphilla group, orange)
'Tom Thumb' (compact, pink sepals, mauve petals)
Pests and diseases
Fuchsias are eaten by the caterpillars of some Lepidoptera, such as the elephant hawk-moth (Deilephila elpenor) and the black-lyre leafroller moth ("Cnephasia" jactatana). Other major insect pests include aphids, mirid bugs such as Lygocoris, Lygus and Plesiocoris spp., vine weevils (Otiorhynchus spp.), and greenhouse whitefly (Trialeurodes vaporariorum). Problematic mites include the fuchsia gall mite (Aculops fuchsiae) and red spider mite (Tetranychus urticae).
Pronunciation and spelling
While the original pronunciation from the word's German origin is , the standard pronunciation for the common name in English is . As a consequence, fuchsia is often misspelled as in English.
When pronounced as scientific Latin name, the pronunciation would be , if one applies the rule that the root word in honorific Latin names should follow as much as possible the original pronunciation of the name of the person the plant is named for, plus the standard pronunciation of the Latin suffix. In practice, however, English-speaking botanists often pronounce it the same as the common name.
History
Leonhart Fuchs, the eminent namesake of the genus, was born in 1501 in Wemding in the Duchy of Bavaria. A physician and professor, he occupied the chair of Medicine at the University of Tübingen from his appointment at the age of 34 until his death in 1566. Besides his medical knowledge, according to his record of activities which was extensive for the time, he studied plants. This was usual for the period. Most remedies and medicines were herbal and the two subjects were often inseparable. In the course of his career Fuchs wrote the seminal De Historia Stirpium Commentarii Insignes, which was richly illustrated and published in 1542. Along with Otto Brunfels (1489–1534) and Hieronymus Bock (1498–1554), also called Hieronymus Tragus, he is today considered one of the three fathers of botany.
It was in honour of Fuchs' and his work that the fuchsia received its name shortly before 1703 by Charles Plumier. Plumier compiled his Nova Plantarum Americanum, which was published in Paris in 1703, based on the results of his third plant-finding trip to the Caribbean in search of new genera. In it he described Fuchsia triphylla flore coccinea.... Plumier's novel species was accepted by Linnaeus in 1753 but the long descriptive name was shortened in accordance with his binomial system.
The first fuchsia species were introduced into English gardens and glasshouses at the end of the 18th century. Fuchsia coccinea Aiton arrived at Kew Gardens in 1788 to be formally described in 1789. It was apparently shortly followed by Fuchsia magellanica Lam. There is much early confusion between these two similar-looking species in the Quelusia Section and they seem to have hybridized readily as well. Fuchsia magellanica, however, proved very hardy outdoors and its cultivars soon naturalized in favorable areas of the British Isles. Other species were quickly introduced to greenhouses. Of special interest is the introduction of Fuchsia fulgens Moç. & Sessé ex DC in the 1830s as it resulted in an outpouring of new cultivars when crossed with the existing species.
Philip Munz, in his A Revision of the Genus Fuchsia (1943), repeats the story that the fuchsia was first introduced into England by a sailor who grew it in a window where it was observed by a nurseryman from Hammersmith, a Mr. Lee, who succeeded in buying it and propagating it for the trade. This was supposedly either one of the short-tubed species such as Fuchsia magellanica or Fuchsia coccinea. The story given by Munz first appears in the 1850s and is embellished in various early publications. Captain Firth, a sailor, brought the plant back to England from one of his trips to his home in Hammersmith where he gave it to his wife. Later James Lee of St. Johns Wood, nurseryman and an astute businessman, heard of the plant and purchased it for £80. He then propagated as many as possible and sold them to the trade for prices ranging from £10 to £20 each. In the Floricultural Cabinet, 1855, there is a report which varies slightly from the above. There it is stated that F. coccinea was given to Kew Garden in 1788 by Captain Firth and that Lee acquired it from Kew. Other than a citation at Kew itself that Fuchsia coccinea was indeed given to it by a Captain Firth, there is no firm evidence to support any of these introduction stories.
Throughout the nineteenth century, plant-collecting fever spread throughout Europe and the United States. Many species of numerous genera were introduced, some as living plants, others as seed. The following fuchsias were recorded in England at Kew: F. lycioides, 1796; F. arborescens, 1824; F. microphylla, 1827; F. fulgens, 1830; F. corymbiflora, 1840; and F. apetala, F. decussata, F. dependens and F. serratifolia in 1843 and 1844, the last four species attributable to Messrs. Veitch of Exeter.
With the increasing numbers of differing species in England plant breeders began to immediately develop hybrids to develop more desirable garden plants. The first recorded experiments date to 1825 as F. arborescens Χ F. macrostemma and F. arborescens X F. coccinea where the quality of the resultant plants was unrecorded.
Between 1835 and 1850 there was a tremendous influx to England of both hybrids and varieties, the majority of which have been lost.
In 1848 Felix Porcher published the second edition of his book Le Fuchsia son Histoire et sa Culture. This described 520 cultivars. In 1871 in later editions of M. Porchers book reference is made to James Lye who was to become famous as a breeder of fuchsias in England. In 1883 the first book of English fuchsias was published.
Between 1900 and 1914 many of the famous cultivated varieties were produced which were grown extensively for Covent Garden market by many growers just outside London. During the period between the world wars, fuchsia-growing slowed as efforts were made toward crop production until after 1949, when plant and hybrid production resumed on a large scale.
In the United States, Sidney Mitchell, a member of the newly formed American Fuchsia Society in San Francisco (1929), shipped a large collection of fuchsias back to California from a nine-month trip to visit gardens in Europe in 1930. Almost immediately after the Society had been established in 1929, a thorough census and collection of fuchsias already growing in California gardens and nurseries had been undertaken under the scientific leadership and direction of Alice Eastwood. The census yielded ninety-one existing cultivars. Armed with that list, Mitchell acquired 51 new fuchsias; 48 of his plants survived the long trip. These were doled out to members of the society and local businesses. Half were also cultivated at the University of California Botanical Garden in Berkeley and the other half at the Berkeley Horticultural Nursery. A wave of interest in fuchsia breeding was launched. Together with the hybrids already in California, many famous American hybrids of the 1940s and 1950s are the descendants of this 1930 group.
| Biology and health sciences | Myrtales | Plants |
268420 | https://en.wikipedia.org/wiki/Foam | Foam | Foams are two-phase material systems where a gas is dispersed in a second, non-gaseous material, specifically, in which gas cells are enclosed by a distinct liquid or solid material. The foam "may contain more or less liquid [or solid] according to circumstances", although in the case of gas-liquid foams, the gas occupies most of the volume. The word derives from the medieval German and otherwise obsolete veim, in reference to the "frothy head forming in the glass once the beer has been freshly poured" (cf. ausgefeimt).
Theories regarding foam formation, structure, and properties—in physics and physical chemistry—differ somewhat between liquid and solid foams in that the former are dynamic (e.g., in their being "continuously deformed"), as a result of gas diffusing between cells, liquid draining from the foam into a bulk liquid, etc. Theories regarding liquid foams have as direct analogs theories regarding emulsions, two-phase material systems in which one liquid is enclosed by another.
In most foams, the volume of gas is large, with thin films of liquid or solid separating the regions of gas. A bath sponge and the head on a glass of beer are examples of foams; soap foams are also known as suds.
Solid foams can be closed-cell or open-cell. In closed-cell foam, the gas forms discrete pockets, each completely surrounded by the solid material. In open-cell foam, gas pockets connect to each other. A bath sponge is an example of an open-cell foam: water easily flows through the entire structure, displacing the air. A sleeping mat is an example of a product composed of closed-cell foam.
Foams are examples of dispersed media. In general, gas is present, so it divides into gas bubbles of different sizes (i.e., the material is polydisperse)—separated by liquid regions that may form films, thinner and thinner when the liquid phase drains out of the system films. When the principal scale is small, i.e., for a very fine foam, this dispersed medium can be considered a type of colloid.
Foam can also refer to something that is analogous to foam, such as quantum foam.
Structure
A foam is, in many cases, a multi-scale system.
One scale is the bubble: material foams are typically disordered and have a variety of bubble sizes. At larger sizes, the study of idealized foams is closely linked to the mathematical problems of minimal surfaces and three-dimensional tessellations, also called honeycombs. The Weaire–Phelan structure is reported in one primary philosophical source to be the best possible (optimal) unit cell of a perfectly ordered foam, while Plateau's laws describe how soap-films form structures in foams.
At lower scale than the bubble is the thickness of the film for metastable foams, which can be considered a network of interconnected films called lamellae. Ideally, the lamellae connect in triads and radiate 120° outward from the connection points, known as Plateau borders.
An even lower scale is the liquid–air interface at the surface of the film. Most of the time this interface is stabilized by a layer of amphiphilic structure, often made of surfactants, particles (Pickering emulsion), or more complex associations.
Formation
Several conditions are needed to produce foam: there must be mechanical work, surface active components (surfactants) that reduce the surface tension, and the formation of foam faster than its breakdown.
To create foam, work (W) is needed to increase the surface area (ΔA):
where γ is the surface tension.
One of the ways foam is created is through dispersion, where a large amount of gas is mixed with a liquid. A more specific method of dispersion involves injecting a gas through a hole in a solid into a liquid. If this process is completed very slowly, then one bubble can be emitted from the orifice at a time as shown in the picture below.
One of the theories for determining the separation time is shown below; however, while this theory produces theoretical data that matches the experimental data, detachment due to capillarity is accepted as a better explanation.
The buoyancy force acts to raise the bubble, which is
where is the volume of the bubble, is the acceleration due to gravity, and ρ1 is the density of the gas ρ2 is the density of the liquid. The force working against the buoyancy force is the surface tension force, which is
,
where γ is the surface tension, and is the radius of the orifice.
As more air is pushed into the bubble, the buoyancy force grows quicker than the surface tension force. Thus, detachment occurs when the buoyancy force is large enough to overcome the surface tension force.
In addition, if the bubble is treated as a sphere with a radius of and the volume is substituted in to the equation above, separation occurs at the moment when
Examining this phenomenon from a capillarity viewpoint for a bubble that is being formed very slowly, it can be assumed that the pressure inside is constant everywhere. The hydrostatic pressure in the liquid is designated by . The change in pressure across the interface from gas to liquid is equal to the capillary pressure; hence,
where R1 and R2 are the radii of curvature and are set as positive. At the stem of the bubble, R3 and R4 are the radii of curvature also treated as positive. Here the hydrostatic pressure in the liquid has to take into account z, the distance from the top to the stem of the bubble. The new hydrostatic pressure at the stem of the bubble is p0(ρ1 − ρ2)z. The hydrostatic pressure balances the capillary pressure, which is shown below:
Finally, the difference in the top and bottom pressure equals the change in hydrostatic pressure:
At the stem of the bubble, the shape of the bubble is nearly cylindrical; consequently, either R3 or R4 is large while the other radius of curvature is small. As the stem of the bubble grows in length, it becomes more unstable as one of the radius grows and the other shrinks. At a certain point, the vertical length of the stem exceeds the circumference of the stem and due to the buoyancy forces the bubble separates and the process repeats.
Stability
Stabilization
The stabilization of foam is caused by van der Waals forces between the molecules in the foam, electrical double layers created by dipolar surfactants, and the Marangoni effect, which acts as a restoring force to the lamellae.
The Marangoni effect depends on the liquid that is foaming being impure. Generally, surfactants in the solution decrease the surface tension. The surfactants also clump together on the surface and form a layer as shown below.
For the Marangoni effect to occur, the foam must be indented as shown in the first picture. This indentation increases the local surface area. Surfactants have a larger diffusion time than the bulk of the solution—so the surfactants are less concentrated in the indentation.
Also, surface stretching makes the surface tension of the indented spot greater than the surrounding area. Consequentially—since the diffusion time for the surfactants is large—the Marangoni effect has time to take place. The difference in surface tension creates a gradient, which instigates fluid flow from areas of lower surface tension to areas of higher surface tension. The second picture shows the film at equilibrium after the Marangoni effect has taken place.
Curing a foam solidifies it, making it indefinitely stable at STP.
Destabilization
Witold Rybczynski and Jacques Hadamard developed an equation to calculate the velocity of bubbles that rise in foam with the assumption that the bubbles are spherical with a radius .
with velocity in units of centimeters per second. ρ1 and ρ2 is the density for a gas and liquid respectively in units of g/cm3 and ῃ1 and ῃ2 is the dynamic
viscosity of the gas and liquid respectively in units of g/cm·s and g is the acceleration of gravity in units of cm/s2.
However, since the density and viscosity of a liquid is much greater than the gas, the density and viscosity of the gas can be neglected, which yields the new equation for velocity of bubbles rising as:
However, through experiments it has been shown that a more accurate model for bubbles rising is:
Deviations are due to the Marangoni effect and capillary pressure, which affect the assumption that the bubbles are spherical.
For laplace pressure of a curved gas liquid interface, the two principal radii of curvature at a point are R1 and R2. With a curved interface, the pressure in one phase is greater than the pressure in another phase. The capillary pressure Pc is given by the equation of:
,
where is the surface tension. The bubble shown below is a gas (phase 1) in a liquid (phase 2) and point A designates the top of the bubble while point B designates the bottom of the bubble.
At the top of the bubble at point A, the pressure in the liquid is assumed to be p0 as well as in the gas. At the bottom of the bubble at point B, the hydrostatic pressure is:
where ρ1 and ρ2 is the density for a gas and liquid respectively. The difference in hydrostatic pressure at the top of the bubble is 0, while the difference in hydrostatic pressure at the bottom of the bubble across the interface is gz(ρ2 − ρ1). Assuming that the radii of curvature at point A are equal and denoted by RA and that the radii of curvature at point B are equal and denoted by RB, then the difference in capillary pressure between point A and point B is:
At equilibrium, the difference in capillary pressure must be balanced by the difference in hydrostatic pressure. Hence,
Since, the density of the gas is less than the density of the liquid the left hand side of the equation is always positive. Therefore, the inverse of RA must be larger than the RB. Meaning that from the top of the bubble to the bottom of the bubble the radius of curvature increases. Therefore, without neglecting gravity the bubbles cannot be spherical. In addition, as z increases, this causes the difference in RA and RB too, which means the bubble deviates more from its shape the larger it grows.
Foam destabilization occurs for several reasons. First, gravitation causes drainage of liquid to the foam base, which Rybczynski and Hadamar include in their theory; however, foam also destabilizes due to osmotic pressure causes drainage from the lamellas to the Plateau borders due to internal concentration differences in the foam, and Laplace pressure causes diffusion of gas from small to large bubbles due to pressure difference. In addition, films can break under disjoining pressure, These effects can lead to rearrangement of the foam structure at scales larger than the bubbles, which may be individual (T1 process) or collective (even of the "avalanche" type).
Mechanical properties
Liquid foams
Solid foams
Solid foams, both open-cell and closed-cell, are considered as a sub-class of cellular structures. They often have lower nodal connectivity as compared to other cellular structures like honeycombs and truss lattices, and thus, their failure mechanism is dominated by bending of members. Low nodal connectivity and the resulting failure mechanism ultimately lead to their lower mechanical strength and stiffness compared to honeycombs and truss lattices.
The strength of foams can be impacted by the density, the material used, and the arrangement of the cellular structure (open vs closed and pore isotropy). To characterize the mechanical properties of foams, compressive stress-strain curves are used to measure their strength and ability to absorb energy since this is an important factor in foam based technologies.
Elastomeric foam
For elastomeric cellular solids, as the foam is compressed, first it behaves elastically as the cell walls bend, then as the cell walls buckle there is yielding and breakdown of the material until finally the cell walls crush together and the material ruptures. This is seen in a stress-strain curve as a steep linear elastic regime, a linear regime with a shallow slope after yielding (plateau stress), and an exponentially increasing regime. The stiffness of the material can be calculated from the linear elastic regime where the modulus for open celled foams can be defined by the equation:
where is the modulus of the solid component, is the modulus of the honeycomb structure, is a constant having a value close to one, is the density of the honeycomb structure, and is the density of the solid. The elastic modulus for closed cell foams can be described similarly by:
where the only difference is the exponent in the density dependence. However, in real materials, a closed-cell foam has more material at the cell edges which makes it more closely follow the equation for open-cell foams. The ratio of the density of the honeycomb structure compared with the solid structure has a large impact on the modulus of the material. Overall, foam strength increases with density of the cell and stiffness of the matrix material.
Energy of deformation
Another important property which can be deduced from the stress strain curve is the energy that the foam is able to absorb. The area under the curve (specified to be before rapid densification at the peak stress), represents the energy in the foam in units of energy per unit volume. The maximum energy stored by the foam prior to rupture is described by the equation:
This equation is derived from assuming an idealized foam with engineering approximations from experimental results. Most energy absorption occurs at the plateau stress region after the steep linear elastic regime.
Directional dependence
The isotropy of the cellular structure and the absorption of fluids can also have an impact on the mechanical properties of a foam. If there is anisotropy present, then the materials response to stress will be directionally dependent, and thus the stress-strain curve, modulus, and energy absorption will vary depending on the direction of applied force. Also, open-cell structures which have connected pores can allow water or other liquids to flow through the structure, which can also affect the rigidity and energy absorption capabilities.
Applications
Liquid foams
Liquid foams can be used in fire retardant foam, such as those that are used in extinguishing fires, especially oil fires.
The dough of leavened bread has traditionally been understood as a closed-cell foam—yeast causing bread to rise via tiny bubbles of gas that become the bread pores—where the cells do not connect with each other. Cutting the dough releases the gas in the bubbles that are cut, but the gas in the rest of the dough cannot escape. When dough is allowed to rise too far, it becomes an open-cell foam, in which the gas pockets are connected; cutting the dough surface at that point would cause a large volume of gas to escape, and the dough to collapse. Recent research has indicated that the pore structure in bread is 99% interconnected into one large vacuole, thus the closed-cell foam of the moist dough is transformed into an open cell solid foam in the bread.
The unique property of gas-liquid foams having very high specific surface area is exploited in the chemical processes of froth flotation and foam fractionation.
Depopulation
Foam depopulation or foaming is a means of mass killing farm animals by spraying foam over a large area to obstruct breathing and ultimately cause suffocation. It is usually used to attempt to stop disease spread.
Solid foams
Solid foams are a class of lightweight cellular engineering materials. These foams are typically classified into two types based on their pore structure: open-cell-structured foams (also known as reticulated foams) and closed-cell foams. At high enough cell resolutions, any type can be treated as continuous or "continuum" materials and are referred to as cellular solids, with predictable mechanical properties.
Open-cell foams can be used to filter air. For example, a foam embedded with catalyst has been shown to catalytically convert formaldehyde to benign substances when formaldehyde polluted air passes through the open cell structure.
Open-cell-structured foams contain pores that are connected to each other and form an interconnected network that is relatively soft. Open-cell foams fill with whatever gas surrounds them. If filled with air, a relatively good insulator results, but, if the open cells fill with water, insulation properties would be reduced. Recent studies have put the focus on studying the properties of open-cell foams as an insulator material. Wheat gluten/TEOS bio-foams have been produced, showing similar insulator properties as for those foams obtained from oil-based resources. Foam rubber is a type of open-cell foam.
Closed-cell foams do not have interconnected pores. The closed-cell foams normally have higher compressive strength due to their structures. However, closed-cell foams are also, in general more dense, require more material, and as a consequence are more expensive to produce. The closed cells can be filled with a specialized gas to provide improved insulation. The closed-cell structure foams have higher dimensional stability, low moisture absorption coefficients, and higher strength compared to open-cell-structured foams. All types of foam are widely used as core material in sandwich-structured composite materials.
The earliest known engineering use of cellular solids is with wood, which in its dry form is a closed-cell foam composed of lignin, cellulose, and air. From the early 20th century, various types of specially manufactured solid foams came into use. The low density of these foams makes them excellent as thermal insulators and flotation devices and their lightness and compressibility make them ideal as packing materials and stuffings.
An example of the use of azodicarbonamide as a blowing agent is found in the manufacture of vinyl (PVC) and EVA-PE foams, where it plays a role in the formation of air bubbles by breaking down into gas at high temperature.
The random or "stochastic" geometry of these foams makes them good for energy absorption, as well. In the late 20th century to early 21st century, new manufacturing techniques have allowed for geometry that results in excellent strength and stiffness per weight. These new materials are typically referred to as engineered cellular solids.
Syntactic foam
Integral skin foam
Integral skin foam, also known as self-skin foam, is a type of foam with a high-density skin and a low-density core. It can be formed in an open-mold process or a closed-mold process. In the open-mold process, two reactive components are mixed and poured into an open mold. The mold is then closed and the mixture is allowed to expand and cure. Examples of items produced using this process include arm rests, baby seats, shoe soles, and mattresses. The closed-mold process, more commonly known as reaction injection molding (RIM), injects the mixed components into a closed mold under high pressures.
Gallery
Foam scales and properties
| Physical sciences | Chemical mixtures: General | null |
268923 | https://en.wikipedia.org/wiki/Elasticity%20%28physics%29 | Elasticity (physics) | In physics and materials science, elasticity is the ability of a body to resist a distorting influence and to return to its original size and shape when that influence or force is removed. Solid objects will deform when adequate loads are applied to them; if the material is elastic, the object will return to its initial shape and size after removal. This is in contrast to plasticity, in which the object fails to do so and instead remains in its deformed state.
The physical reasons for elastic behavior can be quite different for different materials. In metals, the atomic lattice changes size and shape when forces are applied (energy is added to the system). When forces are removed, the lattice goes back to the original lower energy state. For rubbers and other polymers, elasticity is caused by the stretching of polymer chains when forces are applied.
Hooke's law states that the force required to deform elastic objects should be directly proportional to the distance of deformation, regardless of how large that distance becomes. This is known as perfect elasticity, in which a given object will return to its original shape no matter how strongly it is deformed. This is an ideal concept only; most materials which possess elasticity in practice remain purely elastic only up to very small deformations, after which plastic (permanent) deformation occurs.
In engineering, the elasticity of a material is quantified by the elastic modulus such as the Young's modulus, bulk modulus or shear modulus which measure the amount of stress needed to achieve a unit of strain; a higher modulus indicates that the material is harder to deform. The SI unit of this modulus is the pascal (Pa). The material's elastic limit or yield strength is the maximum stress that can arise before the onset of plastic deformation. Its SI unit is also the pascal (Pa).
Overview
When an elastic material is deformed due to an external force, it experiences internal resistance to the deformation and restores it to its original state if the external force is no longer applied. There are various elastic moduli, such as Young's modulus, the shear modulus, and the bulk modulus, all of which are measures of the inherent elastic properties of a material as a resistance to deformation under an applied load. The various moduli apply to different kinds of deformation. For instance, Young's modulus applies to extension/compression of a body, whereas the shear modulus applies to its shear. Young's modulus and shear modulus are only for solids, whereas the bulk modulus is for solids, liquids, and gases.
The elasticity of materials is described by a stress–strain curve, which shows the relation between stress (the average restorative internal force per unit area) and strain (the relative deformation). The curve is generally nonlinear, but it can (by use of a Taylor series) be approximated as linear for sufficiently small deformations (in which higher-order terms are negligible). If the material is isotropic, the linearized stress–strain relationship is called Hooke's law, which is often presumed to apply up to the elastic limit for most metals or crystalline materials whereas nonlinear elasticity is generally required to model large deformations of rubbery materials even in the elastic range. For even higher stresses, materials exhibit plastic behavior, that is, they deform irreversibly and do not return to their original shape after stress is no longer applied. For rubber-like materials such as elastomers, the slope of the stress–strain curve increases with stress, meaning that rubbers progressively become more difficult to stretch, while for most metals, the gradient decreases at very high stresses, meaning that they progressively become easier to stretch. Elasticity is not exhibited only by solids; non-Newtonian fluids, such as viscoelastic fluids, will also exhibit elasticity in certain conditions quantified by the Deborah number. In response to a small, rapidly applied and removed strain, these fluids may deform and then return to their original shape. Under larger strains, or strains applied for longer periods of time, these fluids may start to flow like a viscous liquid.
Because the elasticity of a material is described in terms of a stress–strain relation, it is essential that the terms stress and strain be defined without ambiguity. Typically, two types of relation are considered. The first type deals with materials that are elastic only for small strains. The second deals with materials that are not limited to small strains. Clearly, the second type of relation is more general in the sense that it must include the first type as a special case.
For small strains, the measure of stress that is used is the Cauchy stress while the measure of strain that is used is the infinitesimal strain tensor; the resulting (predicted) material behavior is termed linear elasticity, which (for isotropic media) is called the generalized Hooke's law. Cauchy elastic materials and hypoelastic materials are models that extend Hooke's law to allow for the possibility of large rotations, large distortions, and intrinsic or induced anisotropy.
For more general situations, any of a number of stress measures can be used, and it is generally desired (but not required) that the elastic stress–strain relation be phrased in terms of a finite strain measure that is work conjugate to the selected stress measure, i.e., the time integral of the inner product of the stress measure with the rate of the strain measure should be equal to the change in internal energy for any adiabatic process that remains below the elastic limit.
Units
International System
The SI unit for elasticity and the elastic modulus is the pascal (Pa). This unit is defined as force per unit area, generally a measurement of pressure, which in mechanics corresponds to stress. The pascal and therefore elasticity have the dimension L−1⋅M⋅T−2.
For most commonly used engineering materials, the elastic modulus is on the scale of gigapascals (GPa, 109 Pa).
Linear elasticity
As noted above, for small deformations, most elastic materials such as springs exhibit linear elasticity and can be described by a linear relation between the stress and strain. This relationship is known as Hooke's law. A geometry-dependent version of the idea was first formulated by Robert Hooke in 1675 as a Latin anagram, "ceiiinosssttuv". He published the answer in 1678: "Ut tensio, sic vis" meaning "As the extension, so the force", a linear relationship commonly referred to as Hooke's law. This law can be stated as a relationship between tensile force and corresponding extension displacement ,
where is a constant known as the rate or spring constant. It can also be stated as a relationship between stress and strain :
where is known as the Young's modulus.
Although the general proportionality constant between stress and strain in three dimensions is a 4th-order tensor called stiffness, systems that exhibit symmetry, such as a one-dimensional rod, can often be reduced to applications of Hooke's law.
Finite elasticity
The elastic behavior of objects that undergo finite deformations has been described using a number of models, such as Cauchy elastic material models, Hypoelastic material models, and Hyperelastic material models. The deformation gradient (F) is the primary deformation measure used in finite strain theory.
Cauchy elastic materials
A material is said to be Cauchy-elastic if the Cauchy stress tensor σ is a function of the deformation gradient F alone:
It is generally incorrect to state that Cauchy stress is a function of merely a strain tensor, as such a model lacks crucial information about material rotation needed to produce correct results for an anisotropic medium subjected to vertical extension in comparison to the same extension applied horizontally and then subjected to a 90-degree rotation; both these deformations have the same spatial strain tensors yet must produce different values of the Cauchy stress tensor.
Even though the stress in a Cauchy-elastic material depends only on the state of deformation, the work done by stresses might depend on the path of deformation. Therefore, Cauchy elasticity includes non-conservative "non-hyperelastic" models (in which work of deformation is path dependent) as well as conservative "hyperelastic material" models (for which stress can be derived from a scalar "elastic potential" function).
Hypoelastic materials
A hypoelastic material can be rigorously defined as one that is modeled using a constitutive equation satisfying the following two criteria:
The Cauchy stress at time depends only on the order in which the body has occupied its past configurations, but not on the time rate at which these past configurations were traversed. As a special case, this criterion includes a Cauchy elastic material, for which the current stress depends only on the current configuration rather than the history of past configurations.
There is a tensor-valued function such that in which is the material rate of the Cauchy stress tensor, and is the spatial velocity gradient tensor.
If only these two original criteria are used to define hypoelasticity, then hyperelasticity would be included as a special case, which prompts some constitutive modelers to append a third criterion that specifically requires a hypoelastic model to not be hyperelastic (i.e., hypoelasticity implies that stress is not derivable from an energy potential). If this third criterion is adopted, it follows that a hypoelastic material might admit nonconservative adiabatic loading paths that start and end with the same deformation gradient but do not start and end at the same internal energy.
Note that the second criterion requires only that the function exists. As detailed in the main hypoelastic material article, specific formulations of hypoelastic models typically employ so-called objective rates so that the function exists only implicitly and is typically needed explicitly only for numerical stress updates performed via direct integration of the actual (not objective) stress rate.
Hyperelastic materials
Hyperelastic materials (also called Green elastic materials) are conservative models that are derived from a strain energy density function (W). A model is hyperelastic if and only if it is possible to express the Cauchy stress tensor as a function of the deformation gradient via a relationship of the form
This formulation takes the energy potential (W) as a function of the deformation gradient (). By also requiring satisfaction of material objectivity, the energy potential may be alternatively regarded as a function of the Cauchy-Green deformation tensor (), in which case the hyperelastic model may be written alternatively as
Applications
Linear elasticity is used widely in the design and analysis of structures such as beams, plates and shells, and sandwich composites. This theory is also the basis of much of fracture mechanics.
Hyperelasticity is primarily used to determine the response of elastomer-based objects such as gaskets and of biological materials such as soft tissues and cell membranes.
Factors affecting elasticity
In a given isotropic solid, with known theoretical elasticity for the bulk material in terms of Young's modulus,the effective elasticity will be governed by porosity. Generally a more porous material will exhibit lower stiffness. More specifically, the fraction of pores, their distribution at different sizes and the nature of the fluid with which they are filled give rise to different elastic behaviours in solids.
For isotropic materials containing cracks, the presence of fractures affects the Young and the shear moduli perpendicular to the planes of the cracks, which decrease (Young's modulus faster than the shear modulus) as the fracture density increases, indicating that the presence of cracks makes bodies brittler. Microscopically, the stress–strain relationship of materials is in general governed by the Helmholtz free energy, a thermodynamic quantity. Molecules settle in the configuration which minimizes the free energy, subject to constraints derived from their structure, and, depending on whether the energy or the entropy term dominates the free energy, materials can broadly be classified as energy-elastic and entropy-elastic. As such, microscopic factors affecting the free energy, such as the equilibrium distance between molecules, can affect the elasticity of materials: for instance, in inorganic materials, as the equilibrium distance between molecules at 0 K increases, the bulk modulus decreases. The effect of temperature on elasticity is difficult to isolate, because there are numerous factors affecting it. For instance, the bulk modulus of a material is dependent on the form of its lattice, its behavior under expansion, as well as the vibrations of the molecules, all of which are dependent on temperature.
| Physical sciences | Solid mechanics | null |
269002 | https://en.wikipedia.org/wiki/Aquila%20%28constellation%29 | Aquila (constellation) | Aquila is a constellation on the celestial equator. Its name is Latin for 'eagle' and it represents the bird that carried Zeus/Jupiter's thunderbolts in Greek-Roman mythology.
Its brightest star, Altair, is one vertex of the Summer Triangle asterism. The constellation is best seen in the northern summer, as it is located along the Milky Way. Because of this location, many clusters and nebulae are found within its borders, but they are dim and galaxies are few.
History
Aquila was one of the 48 constellations described by the second-century astronomer Ptolemy. It had been earlier mentioned by Eudoxus in the fourth century BC and Aratus in the third century BC.
It is now one of the 88 constellations defined by the International Astronomical Union. The constellation was also known as Vultur volans (the flying vulture) to the Romans, not to be confused with Vultur cadens which was their name for Lyra. It is often held to represent the eagle which held Zeus's/Jupiter's thunderbolts in Greco-Roman mythology. Aquila is also associated with the eagle that kidnapped Ganymede, a son of one of the kings of Troy (associated with Aquarius), to Mount Olympus to serve as cup-bearer to the gods.
Ptolemy catalogued 19 stars jointly in this constellation and in the now obsolete constellation of Antinous, which was named in the reign of the emperor Hadrian (AD 117–138), but sometimes erroneously attributed to Tycho Brahe, who catalogued 12 stars in Aquila and seven in Antinous. Hevelius determined 23 stars in the first and 19 in the second.
The Greek Aquila is probably based on the Babylonian constellation of the Eagle, but is sometimes mistakenly thought as a seagull which is located in the same area as the Greek constellation.
Notable features
Stars
Aquila, which lies in the Milky Way, contains many rich starfields and has been the location of many novae.
α Aql (Altair) is the brightest star in this constellation and one of the closest naked-eye stars to Earth at a distance of 17 light-years. Its name comes from the Arabic phrase al-nasr al-tair, meaning "the flying eagle". Altair has a magnitude of 0.76. It is one of the three stars of the Summer Triangle, along with Vega and Deneb. It is an A-type main-sequence star with 1.8 times the mass of the Sun and 11 times its luminosity. The star rotates quickly, and this gives the star an oblate shape where it is flattened towards the poles.
β Aql (Alshain) is a yellow-hued star of magnitude 3.7, 45 light-years from Earth. Its name comes from the Arabic phrase shahin-i tarazu, meaning "the balance"; this name referred to Altair, Alshain, and Tarazed. The primary is a G-type subgiant star with a spectral type of G9.5 IV and the secondary is a red dwarf. The subgiant primary has three times the radius of the Sun and six times the luminosity.
γ Aql (Tarazed) is an orange-hued giant star of around magnitude 2.7, 460 light-years from Earth. Its name, like that of Alshain, comes from the Arabic for "the balance". It is the second-brightest star in the constellation and is an unconfirmed variable star.
ζ Aql (Okab) is a binary star of magnitude 3.0, 83 light-years from Earth. The primary is an A-type main sequence star, and the secondary has half the mass of the Sun.
η Aql is a yellow-white-hued supergiant star, 1200 light-years from Earth. Among the brightest Cepheid variable stars, it has a minimum magnitude of 4.4 and a maximum magnitude of 3.5 with a period of 7.2 days. The variability was originally observed by Edward Pigott in 1784. There are also two companion stars which orbit the supergiant: a B-type main sequence star and an F-type main sequence star.
ρ Aql moved across the border into neighboring Delphinus in 1992, and is an A-type star with a lower metallicity than the Sun.
15 Aql is an optical double star. The primary is an orange-hued giant of magnitude 5.41 and a spectral type of K1 III, 325 light-years from Earth. The secondary is a purple-hued star of magnitude 7.0, 550 light-years from Earth. The pair is easily resolved in small amateur telescopes.
57 Aql is a binary star. The primary is a blue-hued star of magnitude 5.7 and the secondary is a white star of magnitude 6.5. The system is approximately 350 light-years from Earth; the pair is easily resolved in small amateur telescopes. Both stars in the system rotate rapidly.
R Aql is a red-hued giant star 690 light-years from Earth. It is a Mira variable with a minimum magnitude of 12.0, a maximum magnitude of 6.0, and a period around 9 months. It has a diameter of 400 D☉.
V Aql is a typical Cool Carbon Star. It's one of the most red-colored examples of this sort of stars, observable through common amateur telescopes.
FF Aql is a yellow-white-hued supergiant star, 2500 light-years from Earth. It is a Cepheid variable star with a minimum magnitude of 5.7, a maximum magnitude of 5.2, and a period of 4.5 days. It is a spectroscopic binary with a spectral type of F6Ib. A third star is also a member of the system, and there is also a fourth star which is probably unconnected with the main system.
Novae
A bright nova was observed in Aquila in 1918 (Nova Aquilae 1918) and briefly shone brighter than Altair, the brightest star in Aquila. It was first seen by Zygmunt Laskowski and was confirmed on the night of 8 June 1918. Nova Aquilae reached a peak apparent magnitude of −0.5 and was the brightest nova recorded since the invention of the telescope.
Deep-sky objects
Three interesting planetary nebulae lie in Aquila:
NGC 6804 shows a small but bright ring.
NGC 6781 bears some resemblance with the Owl Nebula in Ursa Major. It was discovered by William Herschel in 1788.
NGC 6751, also known as the Glowing Eye, is a planetary nebula. The nebula is estimated to be roughly 0.8 light-years in diameter.
More deep-sky objects:
NGC 6709 is a loose open cluster containing roughly 40 stars, which range in magnitude from 9 to 11. It is about 3000 light-years from Earth. It has an overall magnitude of 6.7 and is about 9100 light-years from Earth. NGC 6709 appears in a rich Milky Way star field and is classified as a Shapley class d and Trumpler class III 2 m cluster. These designations mean that it does not have many stars, is loose, does not show greater concentration at the center, and has a moderate range of star magnitudes. There are 305 confirmed member stars and one candidate red giant.
NGC 6755 is an open cluster of 7.5 m; it is made up of about a dozen stars with magnitudes 12 through 13. It is located approximately 8,060 light-years from the Solar System.
NGC 6760 is a globular cluster of 9.1 m. At least two pulsars have been discovered in the globular cluster, and it has a Shapley-Sawyer Concentration Class of IX.
NGC 6749 is an open cluster.
NGC 6778 is a planetary nebula located about 10,300 light-years away from the Solar System.
NGC 6741 is a planetary nebula.
NGC 6772 is a planetary nebula.
W51 (3C400) is one of the largest stellar nurseries in the Milky Way. Located about 17,000 light-years from Earth, W51 is about 350 light-years – or about 2 quadrillion miles – across. However, it's located in an area so thick with interstellar dust that it's opaque to visible light. Observations by the Chandra X-Ray Observatory and the Spitzer Infrared Telescope reveal W51 would appear about as large as the full Moon in visible light.
Aquila also holds some extragalactic objects. One of them is what may be the largest single mass concentration of galaxies in the Universe known, the Hercules–Corona Borealis Great Wall. It was discovered in November 2013, and has the size of 10 billion light years. It is the biggest and the most massive structure in the Universe known.
Other
NASA's Pioneer 11 space probe, which flew by Jupiter and Saturn in the 1970s, is expected to pass near the star Lambda (λ) Aquilae in about 4 million years.
Illustrations
In illustrations of Aquila that represent it as an eagle, a nearly straight line of three stars symbolizes part of the wings. The center and brightest of these three stars is Altair.
Mythology
According to Gavin White, the Babylonian Eagle carried the constellation called the Dead Man in its talons. The author also draws a comparison to the classical stories of Antinous and Ganymede.
In classical Greek mythology, Aquila was identified as Αετός Δίας (Aetos Dios), the eagle that carried the thunderbolts of Zeus and was sent by him to carry the shepherd boy Ganymede, whom he desired, to Mount Olympus; the constellation of Aquarius is sometimes identified with Ganymede.
In the Chinese love story of Qi Xi, Niu Lang (Altair) and his two children (β and γ Aquilae) are separated forever from their wife and mother Zhi Nu (Vega), who is on the far side of the river, the Milky Way.
In Hinduism, the constellation Aquila is identified with the half-eagle half-human deity Garuda.
In ancient Egypt, Aquila possibly was seen as the falcon of Horus. According to Berio, the identification of Aquila as an Egyptian constellation, and not merely Graeco-Babylonian, is corroborated by the Daressy Zodiac. It depicts an outer ring showing the Sphaera Graeca, the familiar Hellenistic zodiac, while the middle ring depicts the Sphaera Barbarica or foreigner's zodiac with the zodiacal signs of the Egyptian dodekaoros which were also recorded by Teucros of Babylon. Under the sign of Sagittarius is the falcon of Horus, presumably because Aquila rises with Sagittarius.
Equivalents
In Chinese astronomy, ζ Aql is located within the Heavenly Market Enclosure (天市垣, Tiān Shì Yuán), and the other stars of the constellation are placed within the Black Tortoise of the North (北方玄武, Běi Fāng Xuán Wǔ).
Several different Polynesian equivalents to Aquila as a whole are known. On the island of Futuna, it was called Kau-amonga, meaning "Suspended Burden". Its name references the Futunan name for Orion's belt and sword, Amonga. In Hawaii, Altair was called Humu, translated to English as "to sew, to bind together parts of a fishhook." "Humu" also refers to the hole by which parts of a hook are bound together. Humu-ma was said to influence the astrologers. Pao-toa was the name for the entire constellation in the Marquesas Islands; the name meant "Fatigued Warrior". Also, Polynesian constellations incorporated the stars of modern Aquila. The Pukapuka constellation Tolu, meaning "three", was made up of Alpha, Beta, and Gamma Aquilae. Altair was commonly named among Polynesian peoples, as well. The people of Hawaii called it Humu, the people of the Tuamotus called it Tukituki ("Pound with a hammer") - they named Beta Aquilae Nga Tangata ("The Men") - and the people of Pukapuka called Altair Turu and used it as a navigational star. The Māori people named Altair Poutu-te-rangi, "Pillar of the Sky", because of its important position in their cosmology. It was used differently in different Māori calendars, being the star of February and March in one version and March and April in the other. Altair was also the star that ruled the annual sweet potato harvest.
| Physical sciences | Other | Astronomy |
269357 | https://en.wikipedia.org/wiki/Wildcat | Wildcat | The wildcat is a species complex comprising two small wild cat species: the European wildcat (Felis silvestris) and the African wildcat (F. lybica). The European wildcat inhabits forests in Europe, Anatolia and the Caucasus, while the African wildcat inhabits semi-arid landscapes and steppes in Africa, the Arabian Peninsula, Central Asia, into western India and western China.
The wildcat species differ in fur pattern, tail, and size: the European wildcat has long fur and a bushy tail with a rounded tip; the smaller African wildcat is more faintly striped, has short sandy-gray fur and a tapering tail; the Asiatic wildcat (F. lybica ornata) is spotted.
The wildcat and the other members of the cat family had a common ancestor about 10–15 million years ago. The European wildcat evolved during the Cromerian Stage about 866,000 to 478,000 years ago; its direct ancestor was Felis lunensis. The silvestris and lybica lineages probably diverged about 173,000 years ago.
The wildcat is categorized as Least Concern on the IUCN Red List since 2002, since it is widely distributed in a stable global population exceeding 20,000 mature individuals. Some local populations are threatened by introgressive hybridisation with the domestic cat (F. catus), contagious disease, vehicle collisions and persecution.
The association of African wildcats and humans appears to have developed along with the establishment of settlements during the Neolithic Revolution, when rodents in grain stores of early farmers attracted wildcats. This association ultimately led to it being tamed and domesticated: the domestic cat is the direct descendant of the African wildcat. It was one of the revered cats in ancient Egypt. The European wildcat has been the subject of mythology and literature.
Taxonomy
Felis (catus) silvestris was the scientific name used in 1777 by Johann von Schreber when he described the European wildcat based on descriptions and names proposed by earlier naturalists such as Mathurin Jacques Brisson, Ulisse Aldrovandi and Conrad Gessner.
Felis lybica was the name proposed in 1780 by Georg Forster, who described an African wildcat from Gafsa on the Barbary Coast.
In subsequent decades, several naturalists and explorers described 40 wildcat specimens collected in European, African and Asian range countries. In the 1940s, the taxonomist Reginald Innes Pocock reviewed the collection of wildcat skins and skulls in the Natural History Museum, London, and designated seven F. silvestris subspecies from Europe to Asia Minor, and 25 F. lybica subspecies from Africa, and West to Central Asia. Pocock differentiated the:
Forest wildcat subspecies (silvestris group)
Steppe wildcat subspecies (ornata-caudata group): is distinguished from the forest wildcat by being smaller, with comparatively lighter fur colour, and longer and more sharply-pointed tails. The domestic cat is thought to have derived from this group.
Bush wildcat subspecies (ornata-lybica group): is distinguished from the steppe wildcat by paler fur, well-developed spot patterns and bands.
In 2005, 22 subspecies were recognized by the authors of Mammal Species of the World, who allocated subspecies largely in line with Pocock's assessment.
In 2017, the Cat Classification Task Force revised the taxonomy of the Felidae, and recognized the following as valid taxa:
Evolution
The wildcat is a member of the Felidae, a family that had a common ancestor about 10–15 million years ago. Felis species diverged from the Felidae around 6–7 million years ago. The European wildcat diverged from Felis about 1.09 to 1.4 million years ago.
The European wildcat's direct ancestor was Felis lunensis, which lived in Europe in the late Pliocene and Villafranchian periods. Fossil remains indicate that the transition from lunensis to silvestris was completed by the Holstein interglacial about 340,000 to 325,000 years ago.
Craniological differences between the European and African wildcats indicate that the wildcat probably migrated during the Late Pleistocene from Europe into the Middle East, giving rise to the steppe wildcat phenotype.
Phylogenetic research revealed that the lybica lineage probably diverged from the silvestris lineage about 173,000 years ago.
Characteristics
The wildcat has pointed ears, which are moderate in length and broad at the base.
Its whiskers are white, number 7 to 16 on each side and reach in length on the muzzle. Whiskers are also present on the inner surface of the paw and measure .
Its eyes are large, with vertical pupils and yellowish-green irises. The eyelashes range from in length, and can number six to eight per side.
The European wildcat has a greater skull volume than the domestic cat, a ratio known as Schauenberg's index. Further, its skull is more spherical in shape than that of the jungle cat (F. chaus) and leopard cat (Prionailurus bengalensis). Its dentition is relatively smaller and weaker than the jungle cat's.
Both wildcat species are larger than the domestic cat. The European wildcat has relatively longer legs and a more robust build compared to the domestic cat. The tail is long, and usually slightly exceeds one-half of the animal's body length. The species size varies according to Bergmann's rule, with the largest specimens occurring in cool, northern areas of Europe and Asia such as Mongolia, Manchuria and Siberia. Males measure in head to body length, in tail length, and normally weigh . Females are slightly smaller, measuring in body length and in tail length, and weighing .
Both sexes have two thoracic and two abdominal teats. Both sexes have pre-anal glands, consisting of moderately sized sweat and sebaceous glands around the anal opening. Large-sized sebaceous and scent glands extend along the full length of the tail on the dorsal side. Male wildcats have pre-anal pockets on the tail, activated upon reaching sexual maturity, play a significant role in reproduction and territorial marking.
Distribution and habitat
The European wildcat inhabits temperate broadleaf and mixed forests in Europe, Turkey and the Caucasus. In the Iberian Peninsula, it occurs from sea level to in the Pyrenees. Between the late 17th and mid 20th centuries, its European range became fragmented due to large-scale hunting and regional extirpation. It is possibly extinct in the Czech Republic, and considered regionally extinct in Austria, though vagrants from Italy are spreading into Austria. It has never inhabited Fennoscandia or Estonia. Sicily is the only island in the Mediterranean Sea with a native wildcat population.
The African wildcat lives in a wide range of habitats except rainforest, but throughout the savannahs of Africa from Mauritania on the Atlantic coast eastward to the Horn of Africa up to altitudes of . Small populations live in the Sahara and Nubian Deserts, Karoo region, Kalahari and Namib Deserts. It occurs around the Arabian Peninsula's periphery to the Caspian Sea, encompassing Mesopotamia, Israel and Palestine region. In Central Asia, it ranges into Xinjiang and southern Mongolia, and in South Asia into the Thar Desert and arid regions in India.
Behaviour and ecology
Both wildcat species are largely nocturnal and solitary, except during the breeding period and when females have young. The size of home ranges of females and males varies according to terrain, the availability of food, habitat quality and the age structure of the population. Male and female home ranges overlap, though core areas within territories are avoided by other cats. Females tend to be more sedentary than males, as they require an exclusive hunting area when raising kittens. Wildcats usually spend the day in a hollow tree, a rock crevice or in dense thickets.
It is also reported to shelter in abandoned burrows of other species such as of red fox (Vulpes vulpes) and in European badger (Meles meles) setts in Europe, and of fennec (Vulpes zerda) in Africa.
When threatened, it retreats into a burrow, rather than climb trees. When taking residence in a tree hollow, it selects one low to the ground. Dens in rocks or burrows are lined with dry grasses and bird feathers. Dens in tree hollows usually contain enough sawdust to make lining unnecessary. If the den becomes infested with fleas, the wildcat shifts to another den. During winter, when snowfall prevents the European wildcat from travelling long distances, it remains within its den until travel conditions improve.
Territorial marking consists of spraying urine on trees, vegetation and rocks, depositing faeces in conspicuous places, and leaving scent marks through glands in its paws. It also leaves visual marks by scratching trees.
Hunting and prey
Sight and hearing are the wildcat's primary senses when hunting.
It lies in wait for prey, then catches it by executing a few leaps, which can span three metres. When hunting near water courses, it waits on trees overhanging the water. It kills small prey by grabbing it in its claws, and piercing the neck or occiput with its fangs. When attacking large prey, it leaps upon the animal's back, and attempts to bite the neck or carotid. It does not persist in attacking if prey manages to escape.
The European wildcat primarily preys on small mammals such as European rabbit (Oryctolagus cuniculus) and rodents.
It also preys on dormice, hares, nutria (Myocastor coypus) and birds, especially ducks and other waterfowl, galliformes, pigeons and passerines. It can consume large bone fragments. Although it kills insectivores such as moles and shrews, it rarely eats them. When living close to human settlements, it preys on poultry. In the wild, it consumes up to of food daily.
The African wildcat preys foremost on murids, to a lesser extent also on birds, small reptiles and invertebrates.
Reproduction and development
The wildcat has two estrus periods, one in December–February and another in May–July. Estrus lasts 5–9 days, with a gestation period lasting 60–68 days. Ovulation is induced through copulation. Spermatogenesis occurs throughout the year. During the mating season, males fight viciously, and may congregate around a single female. There are records of male and female wildcats becoming temporarily monogamous. Kittens are usually born between April and May, and up to August. Litter size ranges from 1–7 kittens.
Kittens are born with closed eyes and are covered in a fuzzy coat. They weigh at birth, and kittens under usually do not survive. They are born with pink paw pads, which blacken at the age of three months, and blue eyes, which turn amber after five months. Their eyes open after 9–12 days, and their incisors erupt after 14–30 days. The kittens' milk teeth are replaced by their permanent dentition at the age of 160–240 days. The kittens start hunting with their mother at the age of 60 days, and start moving independently after 140–150 days. Lactation lasts 3–4 months, though the kittens eat meat as early as 1.5 months of age. Sexual maturity is attained at the age of 300 days. Similarly to the domestic cat, the physical development of African wildcat kittens over the first two weeks of their lives is much faster than that of European wildcats. The kittens are largely fully grown by 10 months, though skeletal growth continues for over 18–19 months. The family dissolves after roughly five months, and the kittens disperse to establish their own territories. Their maximum life span is 21 years, though they usually live up to 13–14 years.
Generation length of the wildcat is about eight years.
Predators and competitors
Because of its habit of living in areas with rocks and tall trees for refuge, dense thickets and abandoned burrows, wildcats have few natural predators. In Central Europe, many kittens are killed by European pine marten (Martes martes), and there is at least one account of an adult wildcat being killed and eaten. Competitors include the golden jackal (Canis aureus), red fox, marten, and other predators. In the steppe regions of Europe and Asia, village dogs constitute serious enemies of wildcats, along with the much larger Eurasian lynx, one of the rare habitual predators of healthy adult wildcats. In Tajikistan, the grey wolf (Canis lupus) is the most serious competitor, having been observed to destroy cat burrows. Birds of prey, including Eurasian eagle-owl (Bubo bubo) and saker falcon (Falco cherrug), have been recorded to kill wildcat kittens. Golden eagle (Aquila chrysaetos) are known to hunt both adults and kittens. Seton Gordon recorded an instance where a wildcat fought a golden eagle, resulting in the deaths of both combatants.
In Africa, wildcats are occasionally killed and eaten by Central African rock python (Python sebae) and martial eagle (Polemaetus bellicosus).
Threats
Wildcat populations are foremost threatened by hybridization with the domestic cat. Mortality due to traffic accidents is a threat especially in Europe. The wildcat population in Scotland has declined since the turn of the 20th century due to habitat loss and persecution by landowners.
In the former Soviet Union, wildcats were caught accidentally in traps set for European pine marten. In modern times, they are caught in unbaited traps on pathways or at abandoned trails of red fox, European badger, European hare or pheasant. One method of catching wildcats consists of using a modified muskrat trap with a spring placed in a concealed pit. A scent trail of pheasant viscera leads the cat to the pit. Wildcat skins were of little commercial value and sometimes converted into imitation sealskin; the fur usually fetched between 50 and 60 kopecks.
Wildcat skins were almost solely used for making cheap scarfs, muffs and coats for ladies.
Conservation
Wildcat species are protected in most range countries and listed in CITES Appendix II. The European wildcat is also listed in Appendix II of the Berne Convention on the Conservation of European Wildlife and Natural Habitats and in the European Union's Habitats and Species Directive.
Conservation Action Plans have been developed in Germany and Scotland.
In culture
Domestication
An African wildcat skeleton excavated in a 9,500-year-old Neolithic grave in Cyprus is the earliest known indication for a close relationship between a human and a possibly tamed cat. As no cat species is native to Cyprus, this discovery indicates that Neolithic farmers may have brought cats to Cyprus from the Near East. Results of genetics and morphological research corroborated that the African wildcat is the ancestor of the domestic cat. The first individuals were probably domesticated in the Fertile Crescent around the time of the introduction of agriculture. Murals and statuettes depicting cats as well mummified cats indicate that it was commonly kept by ancient Egyptians since at least the Twelfth Dynasty of Egypt.
In mythology
Celtic fables of the Cat Sìth, a fairy creature described as resembling a large white-chested black cat, are thought to have been inspired by the Kellas cat, itself thought to be a free-ranging crossbreed between a European wildcat and a domestic cat. In 1693, William Salmon mentioned how body parts of the wildcat were used for medicinal purposes; its flesh for treating gout, its fat for dissolving tumours and easing pain, its blood for curing "falling sickness", and its excrement for treating baldness.
In heraldry
The Picts venerated wildcats, having probably named Caithness (Land of the Cats) after them. According to the foundation myth of the Catti tribe, their ancestors were attacked by wildcats upon landing in Scotland. Their ferocity impressed the Catti so much, that the wildcat became their symbol. The progenitors of Clan Sutherland use the wildcat as symbol on their family crest. The clan's chief bears the title Morair Chat (Great Man of the Cats).
The wildcat is considered an icon of Scottish wilderness, and has been used in clan heraldry since the 13th century. The Clan Chattan Association (also known as the Clan of Cats) comprises 12 clans, the majority of which display the wildcat on their badges.
In literature
Shakespeare referenced the wildcat three times:
| Biology and health sciences | Felines | Animals |
269630 | https://en.wikipedia.org/wiki/Stencil | Stencil | Stencilling produces an image or pattern on a surface by applying pigment to a surface through an intermediate object, with designed holes in the intermediate object. The holes allow the pigment to reach only some parts of the surface creating the design. The stencil is both the resulting image or pattern and the intermediate object; the context in which stencil is used makes clear which meaning is intended. In practice, the (object) stencil is usually a thin sheet of material, such as paper, plastic, wood or metal, with letters or a design cut from it, used to produce the letters or design on an underlying surface by applying pigment through the cut-out holes in the material.
The key advantage of a stencil is that it can be reused to repeatedly and rapidly produce the same letters or design. Although aerosol or painting stencils can be made for one-time use, typically they are made with the intention of being reused. To be reusable, they must remain intact after a design is produced and the stencil is removed from the work surface. With some designs, this is done by connecting stencil islands (sections of material that are inside cut-out "holes" in the stencil) to other parts of the stencil with bridges (narrow sections of material that are not cut out).
Stencil technique in visual art is also referred to as pochoir. A related technique (which has found applicability in some surrealist compositions) is aerography, in which spray-painting is done around a three-dimensional object to create a negative of the object instead of a positive of a stencil design. This technique was used in cave paintings dating to 10,000 BC, where human hands were used in painting handprint outlines among paintings of animals and other objects. The artist sprayed pigment around his hand by using a hollow bone, blown by mouth to direct a stream of pigment.
Screen printing also uses a stencil process, as does mimeography. The masters from which mimeographed pages are printed are often called "stencils". Stencils can be made with one or many colour layers using different techniques, with most stencils designed to be applied as solid colours. During screen printing and mimeography, the images for stenciling are broken down into color layers. Multiple layers of stencils are used on the same surface to produce multi-colored images.
History
Hand stencils, made by blowing pigment over a hand held against a wall, are found from over 35,000 years ago in Asia and Europe, and later prehistoric dates in other continents. After that stenciling has been used as a historic painting technique on all kinds of materials.
Stencils may have been used to color cloth for a very long time; the technique probably reached its peak of sophistication in Katazome and other techniques used on silks for clothes during the Edo period in Japan. In Europe, from about 1450 they were commonly used to color old master prints printed in black and white, usually woodcuts. This was especially the case with playing-cards, which continued to be colored by stencil long after most other subjects for prints were left in black and white. Stencils were used for mass publications, as the type did not have to be hand-written.
Book illustration
Stencils were popular as a method of book illustration, and for that purpose, the technique was at its height of popularity in France during the 1920s when André Marty, Jean Saudé and many other studios in Paris specialized in the technique. Low wages contributed to the popularity of the highly labor-intensive process. When stencils are used in this way they are often called "pochoir".
In the pochoir process, a print with the outlines of the design was produced, and a series of stencils were used through which areas of color were applied by hand to the page. To produce detail, a collotype could be produced which the colors were then stenciled over. Pochoir was frequently used to create prints of intense color and is most often associated with Art Nouveau and Art Deco design.
Aerosol stencils
Aerosol stencils have many practical applications and the stencil concept is used frequently in industrial, commercial, artistic, residential and recreational settings, as well as by the military, government and infrastructure management. A template is used to create an outline of the image. Stencils templates can be made from any material which will hold its form, ranging from plain paper, cardboard, plastic sheets, metals, and wood.
Official use
Stencils are frequently used by official organizations, including the military, utility companies, and governments, to quickly and clearly label objects, vehicles, and locations. Stencils for an official application can be customized, or purchased as individual letters, numbers, and symbols. This allows the user to arrange words, phrases and other labels from one set of templates, unique to the item being labeled. When objects are labeled using a single template alphabet, it makes it easier to identify their affiliation or source.
Stencil graffiti
Stencils have also become popular for graffiti, since stencil art using spray-paint can be produced quickly and easily. These qualities are important for graffiti artists where graffiti is illegal or quasi-legal, depending on the city and stenciling surface. The extensive lettering possible with stencils makes it especially attractive to political artists. For example, the anarcho-punk band Crass used stencils of anti-war, anarchist, feminist and anti-consumerist messages in a long-term graffiti campaign around the London Underground system and on advertising billboards. There has been a semi-recent trend in making multi-layered stencils with different shades of grey for each layer creating a more detailed stenciled image. Also well known for their use of stencil art are Blek le Rat, Epsylon, Marie Rouffet, Nuklé-art, Kim Prisu, Miss Tic and Jef aerosol from France, British artist Banksy, New York artist, world traveling artist Tavar Zawacki f.k.a. 'ABOVE', Shepard Fairey's OBEY, and Pirate & Acid from Hollywood, California.
Home stenciling
A common tradition for stencils is in home decorating and arts & crafts. Home decor stencils are an important part of the DIY (Do It Yourself) industry. There are prefabricated stencil templates available for home decoration projects from hardware stores, arts & crafts stores and through the internet. Stencils are usually applied in the home with a paint or roller brush along wall borders and as trim. They can also be applied with a painted sponge for a textured effect.
Stencil templates can be purchased or constructed individually. Typically they are constructed of flexible plastics, including acetate, mylar, and vinyl. Stencils can be used as children's toys.
Military stenciling
Stencils have been used in the military across most nations for many years and continue to be used today. They are used to mark up equipment, vehicles, rations, signposts, helmets, etc. One use of military stencils was the application of playing card designs to USA Airborne helmets during World War Two as a method to identify regimental units.
Silk screening
Silk screening is a type of printing on paper or textiles, in which an ink is embedded in the cloth. The ink is controlled through the use of a stencil, which is placed directly over the paper or textile. This process can only handle one color of ink at a time. Therefore, multi-colored designs must be silk screened several times, with each interval taking time to dry.
Micro- and nanostencil
Stencils are also used in micro- and nanotechnology, as miniature shadow masks through which material can be deposited, etched or ions implanted onto a substrate. These stencils are usually made out of thin (100-500 nm) low-stress Silicon nitride (SiN) in which apertures are defined by various lithographic techniques (e. g. electron beam, photolithography).
Stencil lithography has unique advantages compared to other patterning techniques: it does not require spinning of a uniform layer of resist (therefore patterns can be created on 3D topographies) and it does not involve any heat or chemical treatment of the substrate (like baking, developing and removing the resist). Thus it allows a wide range of substrates (e.g. flexible, surface-treated) and materials (e. g. organics) to be used.
Other stencil forms
Screen printing
A stencil technique is employed in screen printing which uses a tightly woven mesh screen coated in a thin layer of emulsion to reproduce the original image. As the stencil is attached to the screen, a contiguous template is not necessary.
Airbrushing
A stencil used in airbrushing called a frisket is pressed directly on the artwork. It can be used to control or contain overspray, create sharp or complex shapes, but is not designed to be used more than once.
Wall stencils
Wall stencils - to decorate walls and ceilings or create your own repeat for an overall modern wall pattern effect.
Rock art
One form of pictograph found in ancient and traditional rock paintings is created by the hand first being placed against the panel, with dry paint then being blown onto it through a tube, in a process that is akin to air-brush or spray-painting. The resulting image is a negative print of the hand, and is sometimes described as a "stencil" in Australian archaeology.
Miniature rock art of the stencilled variety at a rock shelter known as Yilbilinji, in the Limmen National Park in the Northern Territory, is one of only three known examples of such art. Usually stencilled art is life-size, using body parts as the stencil, but the 17 images of designs of human figures, boomerangs, animals such as crabs and long-necked turtles, wavy lines and geometric shapes are very rare. Found in 2017 by archaeologists, the only other recorded examples are at Nielson's Creek in New South Wales and at Kisar Island in Indonesia. It is thought that the designs may have been created by stencils fashioned out of beeswax.
| Technology | Printing | null |
269676 | https://en.wikipedia.org/wiki/Wet%20nurse | Wet nurse | A wet nurse is a woman who breastfeeds and cares for another's child. Wet nurses are employed if the mother dies, if she is unable to nurse the child herself sufficiently or chooses not to do so. Wet-nursed children may be known as "milk-siblings", and in some societies, the families are linked by a special relationship of milk kinship. Wet-nursing existed in societies around the world until the invention of reliable formula milk in the 20th century. The practice has made a small comeback in the 21st century.
Reasons
A wet nurse can help when a mother is unable or unwilling to breastfeed her baby. Before the development of infant formula in the 20th century, wet-nursing could save a baby's life.
There are many reasons why a mother is unable to produce sufficient breast milk, or in some cases to lactate at all. For example, she may have a chronic or acute illness, and either the illness itself, or the treatment for it, reduces or stops her milk. This absence of lactation may be temporary or permanent.
There was a greater need for wet nurses when the rates of infant abandonment and maternal death, during and shortly after childbirth, were high. There was a concurrent availability of lactating women whose own babies had died.
Some women chose not to breastfeed for social reasons. For upper-class women, breastfeeding was considered unfashionable, in the sense that it not only prevented them from being able to wear the fashionable clothing of their time, but it was also thought to ruin their figures. Hiring a wet nurse was less expensive than having to hire someone else to help run the family business and/or take care of the family household duties in their place. Some women chose to hire wet nurses purely to escape from the confining and time-consuming chore of breastfeeding.
Eliciting milk
A woman can only act as a wet nurse if she is lactating (producing milk). It was once believed that a wet nurse must have recently undergone childbirth in order to lactate. This is not necessarily the case, as regular breast stimulation can elicit lactation via a neural reflex of prolactin production and secretion. Some women have been able to establish lactation using a breast pump, in order to feed an infant.
Gabrielle Palmer, author of The Politics of Breastfeeding, states:
Historical and cultural practices
Wet nursing is an ancient practice, common to many societies. It has been linked to social class, where monarchies, the aristocracy, nobility, or upper classes had their children wet-nursed for the benefit of the child's health, and sometimes in the hope of becoming pregnant again quickly. Exclusive breastfeeding inhibits ovulation in some women (lactational amenorrhea). Poor women, especially those who suffered the stigma of giving birth to an illegitimate child, sometimes had to give their baby up temporarily to a wet nurse, or permanently to another family. The woman herself might in turn become wet nurse to a wealthier family, while using part of her wages to pay her own child's wet nurse.
In pre-modern times, it was incorrectly believed that wet nurses could pass on personality traits to infants, such as acquired characteristics.
Mythology
Many cultures feature stories, historical or mythological, involving superhuman, supernatural, human, and in some instances, animal wet nurses.
The Bible refers to Deborah, a nurse to Rebekah, wife of Isaac and mother of Jacob (Israel) and Esau, who appears to have lived as a member of the household all her days. (Genesis 35:8.) Midrashic commentaries on the Torah hold that the Egyptian princess Bithiah (Pharaoh's wife Asiya in the Islamic Hadith and Qur'an) attempted to wet-nurse Moses, but he would take only his biological mother's milk. ()
In Greek mythology, Eurycleia is the wet nurse of Odysseus. In Roman mythology, Caieta was the wet nurse of Aeneas. In Burmese mythology, Myaukhpet Shinma is the nat (spirit) representation of the wet nurse of King Tabinshwehti. In Hawaiian mythology, Nuakea is a beneficent goddess of lactation; her name became the title for a royal wet nurse, according to David Malo.
Ancient Rome
In ancient Rome, well-to-do households would have had wet nurses (Latin , singular ) among their slaves and freedwomen, but some Roman women were wet nurses by profession, and the Digest of Roman law even refers to a wage dispute for wet-nursing services (). The landmark known as the Columna Lactaria ("Milk Column") may have been a place where wet nurses could be hired. It was considered admirable for upperclass women to breastfeed their own children, but unusual and old-fashioned in the Imperial era. Even women of the working classes or slaves might have their babies nursed, and the Roman-era Greek gynecologist Soranus offers detailed advice on how to choose a wet nurse. Inscriptions such as religious dedications and epitaphs indicate that a would be proud of her profession. One even records a , a male "milk nurse" who presumably used a bottle. Greek nurses were preferred, and the Romans believed that a baby who had a Greek could imbibe the language and grow up speaking Greek as fluently as Latin.
The importance of the wet nurse to ancient Roman culture is indicated by the founding myth of Romulus and Remus, who were abandoned as infants but nursed by the she-wolf, as portrayed in the famous Capitoline Wolf bronze sculpture. The goddess Rumina was invoked among other birth and child development deities to promote the flow of breast milk.
India
By the 1500s, a wealthy mother who did not use a wet nurse was worthy of remark in India. The child was not "put out" of the household; rather, the wet nurse was included within it. The imperial wet nurses of the Mughal court were given honours in the Turco-Mongol tradition.
United Kingdom
Wet nursing used to be commonplace in the United Kingdom. Working-class women both provided and received wet-nursing services.
Taking care of babies was a well-paid, respectable, and popular job for many working-class women. In the 18th century, a woman would earn more money as a wet nurse than an average man could as a labourer. Up until the 19th century, most wet-nursed infants were sent far from their families to live with their new caregiver for up to the first three years of their life. As many as 80% of wet-nursed babies who lived like this died during infancy.
During the Victorian era, women took in babies for money and nursed them themselves or fed them with whatever was cheapest. This was known as baby-farming; poor care sometimes resulted in high infant death rates. The wet nurse at this period was most likely a single woman who previously had given birth to an illegitimate child. There were two types of wet nurses by this time: those on poor relief, who struggled to provide sufficiently for themselves or their charges, and the professionals, who were well paid and respected.
Upper-class women tended to hire wet nurses to work within their own homes, as part of a large household of servants.
Wet nurses also worked at foundling hospitals, establishments for abandoned children. Their own children would likely be sent away, normally brought up by the bottle rather than being breastfed. Valerie Fildes, author of Breasts, Bottle and Babies: A History of Infant Feeding, argues that "In effect, wealthy parents frequently 'bought' the life of their infant for the life of another."
Wet nursing decreased in popularity during the mid-19th century, as medical journalists wrote about its previously undocumented dangers. Fildes argued that "Britain has been lumped together with the rest of Europe in any discussion of the qualities, terms of employment and conditions of the wet nurse, and particularly the abuses of which she was supposedly guilty." C. H. F. Routh, a medical journalist writing in the late 1850s, listed the evils of wet nursing, such as the abandonment of the wet nurses' own children, higher infant mortality, and an increased physical and moral risk to a nursed child. While this argument was not founded in any sort of proof, the emotional arguments of medical researchers, coupled with the protests of other critics, slowly increased public knowledge; the practice declined, replaced by maternal breastfeeding and bottle-feeding.
France
Wet-nursing was reported in France in the time of Louis XIV, the mid-17th century. By the 18th century, approximately 90% of infants were wet-nursed, mostly sent away to live with their wet nurses. In Paris, only 1,000 of the 21,000 babies born in 1780 were nursed by their own mothers. The high demand for wet nurses coincided with the low wages and high rent prices of this era, which forced many women to have to work soon after childbirth. This meant that many mothers had to send their infants away to be breastfed and cared for by wet nurses even poorer than themselves. With the high demand for wet nurses, the price to hire one increased as the standard of care decreased. This led to many infant deaths. In response, rather than nursing their own children, upper-class women turned to hiring wet nurses to come live with them instead. In entering into their employer's home to care for their charges, these wet nurses had to leave their own infants to be nursed and cared for by women far worse off than themselves, and who likely lived at a relatively far distance away.
The Bureau of Wet Nurses was created in Paris in 1769 to serve two main purposes: it supplied parents with wet nurses, as well as helping lessen the neglect of babies by controlling monthly salary payments. In order to become a wet nurse, women had to meet a few qualifications, including physical fitness and good moral character; they were often judged on their age, their health, the number of children they had, as well as their breast shape, breast size, breast texture, nipple shape, and nipple size, since all these aspects were believed to affect the quality of a woman's milk. In 1874, the French government introduced a law named after , which "mandated that every infant placed with a paid guardian outside the parents' home be registered with the state so that the French government is able to monitor how many children are placed with wet nurses and how many wet-nursed children have died".
Wet nurses were hired to work in hospitals to nurse babies who were premature, ill, or abandoned. During the 18th and 19th centuries, congenital syphilis was a common cause of infant mortality. The Vaugirard hospital in Paris began to use mercury as a treatment; however, it could not be safely administered to infants. In 1780, it began the process of giving mercury to wet nurses, who could then transmit the treatment in their milk to infected infants.
The practice of wet-nursing was still widespread during World War I, according to the American Red Cross. Working-class women would leave their babies with wet nurses so they could get jobs in factories.
United States
British colonists brought the practice of wet-nursing with them to North America. Since the arrangement of sending infants away to live with wet nurses was the cause of so many infant deaths, by the 19th century, Americans adopted the practice of having wet nurses live with the employers in order to nurse and care for their charges. This practice had the effect of increasing the death rate for wet nurses' own babies. Many employers would have only kept a wet nurse for a few months at a time since it was believed that the quality of a woman's breast milk would lessen over time.
Child-minding, different from wet-nursing, was also commonly an additional job on top of child rearing and nursery tending. Employed wet nurses were typically paid low wages and worked long hours. Workers in the 1900s demanded work contracts to provide stable wages. Wet nursing work was rarely consistent, wet nurses were stereotypically poor ladies from rural areas who offered their services for fees.
Since there were no official records kept pertaining to wet nurses or wet-nursed babies, historians lack the knowledge of precisely how many infants were wet-nursed and for how long, whether they lived at home or elsewhere, and how many lived or died. The best source of evidence is found in the "help wanted" ads of newspapers, through complaints about wet nurses in magazines, and through medical journals that acted as employment agencies.
Slavery
In the Southern United States before the Civil War, it was common practice for enslaved black women to be forced to be wet nurses to their owners' children. In some instances, the enslaved child and the white child would be raised together in their younger years. (Sometimes both babies would be fathered by the same man, the slave-owner; see Children of the plantation.) Visual representations of wet-nursing practices in enslaved communities are most prevalent in representations of the Mammy archetype caricature. Images such as the one in this section represent both a historically accurate practice of enslaved black women wet-nursing their owner's white children, as well as sometimes an exaggerated racist caricaturization of a stereotype of a "Mammy" character.
Egypt
From the mid-1800s to the mid-1900s, and especially after World War I, thousands of Slovene peasant women migrated via Trieste to the cosmopolitan port city of Alexandria. There, these undertook various sorts of domestic work for elite Levantine households—"the highly mobile upper strata of Ottoman millets, Jewish, Maronites, Melkite active in international commerce". Enough served as wet nurses that this occupation became almost synonymous with Slovene domestic workers, which resulted in some stigma back home. Married women could leave Alexandria and return to their home village, where they would conceive and bear a child and leave the infant to the care of relatives or a hired wet nurse, while they returned to Egypt to seek new employment and a new charge to nurse.
This constitutes the origin of the archetype of the as a wet nurse, which came to overpower any other representation of the , despite the fact that empirical evidence demonstrates that only a tiny fraction of at any time worked as wet nurses. The majority of were working as nannies or chamber maids, they were not breastfeeding the children they were taking care of. The emphasis on lactaction, which marks the hypersexualization of the , was part of the rhetorical stigma surrounding this phenomenon in Slovenia.
Relationships
Sometimes, the infant was placed in the home of the wet nurse for several months, as was the case for Jane Austen and her siblings. The Papyrus Oxyrhynchus 91, a receipt from AD 187, attests to the ancient nature of this practice. Sometimes, the wet nurse came to live with the infant's family, filling a position between the monthly nurse (for the immediate post-partum period) and the nanny.
In some societies, the wet nurse was simply hired as any other employee. In others, however, she had a special relationship with the family, which could incur kinship rights. In Vietnamese family structure, for example, the wet nurse is known as , meaning "mother". Islam has a highly codified system of milk kinship known as rada. George III of the United Kingdom, born two months premature, had a wet nurse whom he so valued all his life, that her daughter was appointed laundress to the Royal Household, "a sinecure place of great emolument".
Mothers who nurse each other's babies are engaging in a reciprocal act known as cross-nursing or co-nursing.
Current attitudes in Western countries
In contemporary affluent Western societies such as the United States, the act of nursing a baby other than one's own often provokes cultural discomfort. When a mother is unable to nurse her own infant, an acceptable mediated substitute is expressed milk (or especially colostrum), which is donated to milk banks, analogous to blood banks, and processed there by being screened, pasteurized, and usually frozen. Infant formula is also widely available, which its makers claim can be a reliable source of infant nutrition when prepared properly. Dr. Rhonda Shaw notes that Western objections to wet nurses are cultural:
For some Americans, the subject of wet-nursing is becoming increasingly open for discussion. During a UNICEF goodwill tour to Sierra Leone in 2008, American Mexican actress Salma Hayek decided to breastfeed a local infant in front of the accompanying film crew. The sick one-week-old baby had been born the same day but a year later than Hayek's daughter, who had not yet been weaned. The actress later discussed on camera an anecdote of her Mexican great-grandmother spontaneously breastfeeding a hungry baby in a village.
Current situation elsewhere
Wet nurses are still common in many developing countries, although the practice poses a risk of infections, such as HIV. In China, Indonesia, and the Philippines, a wet nurse may be employed in addition to a nanny as a mark of aristocracy, wealth, and high status. Following the 2008 Chinese milk scandal, in which contaminated infant formula poisoned thousands of babies, the salaries of wet nurses there increased dramatically.
Notable wet nurses
Royal wet nurses are more likely than most to reach the historical record.
In Ancient Egypt, Maia was the wet nurse of King Tutankhamun. Sitre In, the nurse of Hatshepsut, was not a member of the royal family but received the honour of a burial in the royal necropolis in the Valley of the Kings, in tomb KV60. Her coffin has the inscription , meaning Great Royal Wet Nurse In.
In Asia, Lady Kasuga was the wet nurse of the third Tokugawa shōgun, Iemitsu. Lu Lingxuan was a lady in waiting who served as wet nurse to the emperor Gao Wei. She became exceedingly powerful during his reign and was often criticized by historians for her corruption and treachery. Chinese emperors honoured the Nurse empress dowager. Wet nurses were also common during the Mughal period, with almost every Mughal prince having one. Some prominent ones are Maham Anga for Akbar and Dai Anga for Shah Jahan. Shin Myo Myat was the mother of King Bayinnaung of the Toungoo dynasty of Burma (Myanmar), and the wet nurse of King Tabinshwehti. The last Emperor of China, Puyi, described Wang Lianshou as being the only person who was able to control him: "from my infancy until the time I grew up, only my wet nurse, because of her simple language, was able to make me grasp the idea that I was like other people."
In Europe, Hodierna of St Albans was the mother of Alexander Neckam and wet nurse of Richard I of England, and Mrs. Pack was a wet nurse to William, Duke of Gloucester (1689–1700). Geneviève Poitrine was a wet nurse of the Dauphin of France, Louis Joseph, son of King Louis XVI of France and Queen Marie Antoinette. Poitrine was accused of transmitting tuberculosis to the Dauphin and triggering his infant death when aged seven, although since very few pre-adolescent children die from TB, this accusation may have been the result of a misdiagnosis.
Some non-royal wet nurses have also been written about. Halimah bint Abi Dhuayb was the foster mother and wet nurse of the Islamic prophet Muhammad. Petronella Muns was, with her employer, the first Western woman to visit Japan. Naomi Baumslag, author of Milk, Money and Madness, described the legendary capacity of Judith Waterford: "In 1831, on her 81st birthday, she could still produce breast milk. In her prime she unfailingly produced two quarts (four pints or 1.9 litres) of breast milk a day."
| Biology and health sciences | Health and fitness: General | Health |
270054 | https://en.wikipedia.org/wiki/Formal%20verification | Formal verification | In the context of hardware and software systems, formal verification is the act of proving or disproving the correctness of a system with respect to a certain formal specification or property, using formal methods of mathematics.
Formal verification is a key incentive for formal specification of systems, and is at the core of formal methods.
It represents an important dimension of analysis and verification in electronic design automation and is one approach to software verification. The use of formal verification enables the highest Evaluation Assurance Level (EAL7) in the framework of common criteria for computer security certification.
Formal verification can be helpful in proving the correctness of systems such as: cryptographic protocols, combinational circuits, digital circuits with internal memory, and software expressed as source code in a programming language. Prominent examples of verified software systems include the CompCert verified C compiler and the seL4 high-assurance operating system kernel.
The verification of these systems is done by ensuring the existence of a formal proof of a mathematical model of the system. Examples of mathematical objects used to model systems are: finite-state machines, labelled transition systems, Horn clauses, Petri nets, vector addition systems, timed automata, hybrid automata, process algebra, formal semantics of programming languages such as operational semantics, denotational semantics, axiomatic semantics and Hoare logic.
Approaches
Model Checking
Model checking involves a systematic and exhaustive exploration of the mathematical model. Such exploration is possible for finite models, but also for some infinite models, where infinite sets of states can be effectively represented finitely by using abstraction or taking advantage of symmetry. Usually, this consists of exploring all states and transitions in the model, by using smart and domain-specific abstraction techniques to consider whole groups of states in a single operation and reduce computing time. Implementation techniques include state space enumeration, symbolic state space enumeration, abstract interpretation, symbolic simulation, abstraction refinement. The properties to be verified are often described in temporal logics, such as linear temporal logic (LTL), Property Specification Language (PSL), SystemVerilog Assertions (SVA), or computational tree logic (CTL). The great advantage of model checking is that it is often fully automatic; its primary disadvantage is that it does not in general scale to large systems; symbolic models are typically limited to a few hundred bits of state, while explicit state enumeration requires the state space being explored to be relatively small.
Deductive Verification
Another approach is deductive verification. It consists of generating from the system and its specifications (and possibly other annotations) a collection of mathematical proof obligations, the truth of which imply conformance of the system to its specification, and discharging these obligations using either proof assistants (interactive theorem provers) (such as HOL, ACL2, Isabelle, Coq or PVS), or automatic theorem provers, including in particular satisfiability modulo theories (SMT) solvers. This approach has the disadvantage that it may require the user to understand in detail why the system works correctly, and to convey this information to the verification system, either in the form of a sequence of theorems to be proved or in the form of specifications (invariants, preconditions, postconditions) of system components (e.g. functions or procedures) and perhaps subcomponents (such as loops or data structures).
Application to Software
Formal verification of software programs involves proving that a program satisfies a formal specification of its behavior. Subareas of formal verification include deductive verification (see above), abstract interpretation, automated theorem proving, type systems, and lightweight formal methods. A promising type-based verification approach is dependently typed programming, in which the types of functions include (at least part of) those functions' specifications, and type-checking the code establishes its correctness against those specifications. Fully featured dependently typed languages support deductive verification as a special case.
Another complementary approach is program derivation, in which efficient code is produced from functional specifications by a series of correctness-preserving steps. An example of this approach is the Bird–Meertens formalism, and this approach can be seen as another form of program synthesis.
These techniques can be sound, meaning that the verified properties can be logically deduced from the semantics, or unsound, meaning that there is no such guarantee. A sound technique yields a result only once it has covered the entire space of possibilities. An example of an unsound technique is one that covers only a subset of the possibilities, for instance only integers up to a certain number, and give a "good-enough" result. Techniques can also be decidable, meaning that their algorithmic implementations are guaranteed to terminate with an answer, or undecidable, meaning that they may never terminate. By bounding the scope of possibilities, unsound techniques that are decidable might be able to be constructed when no decidable sound techniques are available.
Verification and validation
Verification is one aspect of testing a product's fitness for purpose. Validation is the complementary aspect. Often one refers to the overall checking process as V & V.
Validation: "Are we trying to make the right thing?", i.e., is the product specified to the user's actual needs?
Verification: "Have we made what we were trying to make?", i.e., does the product conform to the specifications?
The verification process consists of static/structural and dynamic/behavioral aspects. E.g., for a software product one can inspect the source code (static) and run against specific test cases (dynamic). Validation usually can be done only dynamically, i.e., the product is tested by putting it through typical and atypical usages ("Does it satisfactorily meet all use cases?").
Automated program repair
Program repair is performed with respect to an oracle, encompassing the desired functionality of the program which is used for validation of the generated fix. A simple example is a test-suite—the input/output pairs specify the functionality of the program. A variety of techniques are employed, most notably using satisfiability modulo theories (SMT) solvers, and genetic programming, using evolutionary computing to generate and evaluate possible candidates for fixes. The former method is deterministic, while the latter is randomized.
Program repair combines techniques from formal verification and program synthesis. Fault-localization techniques in formal verification are used to compute program points which might be possible bug-locations, which can be targeted by the synthesis modules. Repair systems often focus on a small pre-defined class of bugs in order to reduce the search space. Industrial use is limited owing to the computational cost of existing techniques.
Industry use
The growth in complexity of designs increases the importance of formal verification techniques in the hardware industry. At present, formal verification is used by most or all leading hardware companies, but its use in the software industry is still languishing. This could be attributed to the greater need in the hardware industry, where errors have greater commercial significance. Because of the potential subtle interactions between components, it is increasingly difficult to exercise a realistic set of possibilities by simulation. Important aspects of hardware design are amenable to automated proof methods, making formal verification easier to introduce and more productive.
, several operating systems have been formally verified:
NICTA's Secure Embedded L4 microkernel, sold commercially as seL4 by OK Labs; OSEK/VDX based real-time operating system ORIENTAIS by East China Normal University; Green Hills Software's Integrity operating system; and SYSGO's PikeOS.
In 2016, a team led by Zhong Shao at Yale developed a formally verified operating system kernel called CertiKOS.
As of 2017, formal verification has been applied to the design of large computer networks through a mathematical model of the network, and as part of a new network technology category, intent-based networking. Network software vendors that offer formal verification solutions include Cisco Forward Networks and Veriflow Systems.
The SPARK programming language provides a toolset which enables software development with formal verification and is used in several high-integrity systems.
The CompCert C compiler is a formally verified C compiler implementing the majority of ISO C.
| Technology | Software development: General | null |
270445 | https://en.wikipedia.org/wiki/Human%20hair%20color | Human hair color | Human hair color is the pigmentation of human hair follicles and shafts due to two types of melanin: eumelanin and pheomelanin. Generally, the more melanin present, the darker the hair. Its tone depends on the ratio of black or brown eumelanin to yellow or red pheomelanin. Melanin levels can vary over time, causing a person's hair color to change, and one person can have hair follicles of more than one color. Some hair colors are associated with some ethnic groups because of the observed higher frequency of particular hair colors within their geographical region, e.g. straight, dark hair amongst East Asians, Southeast Asians, Polynesians, some Central Asians, and Native Americans; a large variety of dark, fair, curly, straight, wavy or bushy amongst Europeans, West Asians, some Central Asians, and North Africans; and curly, dark, and uniquely helical hair amongst Sub Saharan Africans. Bright red hair is found in some European populations, and hair turns gray, white, or "silver" with age.
Genetics and biochemistry of hair color
The full genetic basis of hair color is complex and not fully understood. Regulatory DNA is believed to be closely involved in pigmentation in humans in general, and a 2011 study by Branicki et al. identified 13 DNA variations across 11 different genes that could be used to predict hair color.
Two types of pigment give hair its color, black-brown eumelanin and reddish-brown/reddish-yellow pheomelanin, synthesized by melanocytes. Inside the melanocytes, tyrosine is converted into L-DOPA and then L-dopaquinone, which in turn is formed into pheomelanin or eumelanin.
Different hair color phenotypes arise primarily as a result of varying ratios of these two pigments in the human population, although Europeans show the greatest range in pigmentation overall. In addition, other genetic and environmental factors can affect hair color in humans; for instance, mutations in the melanocortin 1 receptor (MC1R) gene can lead to red or auburn hair, and exposure to ultraviolet radiation can damage hair and alter its pigmentation. Ultraviolet radiation (UV radiation) triggers greater synthesis of several compounds, including pro-opiomelanocortin (POMC), α-MSH, and ACTH, the result being increased eumelanin production. UV radiation most commonly comes from the sun, and thus populations from places closer to the equator tend to have darker hair, because eumelanin is generally more photoprotective than pheomelanin.
Pheomelanin colors hair orange and red. Eumelanin, which has two subtypes of black or brown, determines the darkness of the hair color; more black eumelanin leads to blacker hair and more brown eumelanin to browner hair. All human hair has some amount of both pigments. Over 95% of melanin content in black and brown hair is eumelanin. Pheomelanin is generally found in elevated concentrations in blond and red hair, representing about one-third of total melanin content. If there is no black eumelanin, the result is strawberry blond. Blond hair results from small amounts of brown eumelanin with no black eumelanin.
Natural hair colors
Natural hair color can be black, brown, blonde and red.
Color shade scale
The Fischer–Saller scale, named after Eugen Fischer and Karl Saller is used in physical anthropology and medicine to determine the shades of hair color. The scale uses the following designations: A (very light blond), B to E (light blond), F to L (blond), M to O (dark blond), P to T (light brown to brown), U to Y (dark brown to black) and Roman numerals I to IV (red) and V to VI (red-blond).
Image gallery
Black hair
Black hair or jet black hair is the darkest hair color. It has large amounts of eumelanin and is denser than other hair colors and is the commonly seen hair color in Asia and Africa due the fact that the people in these regions tend to have lower levels of tyrosinase in their bodies. Black eumelanin secretion causes the hair to turn black, which indicates that the MC1R is in the active state. Jet black hair, the darkest shade will not have a warm, neutral tone but a sheen which can seem almost blue, like the iridescence of a raven's wing; hence, sometimes referred to as raven-black. Jet black hair appears to have reflective silver color in bright sunlight.
Brown hair
Brown hair is the second most common human hair color, after black. Brown hair is characterized by higher levels of eumelanin and lower levels of pheomelanin. Of the two types of eumelanin (black and brown), brown-haired people have brown eumelanin; they also usually have medium-thick strands of hair. Brown-haired girls or women of European, West Asian or North African descent are often known as brunettes.
Chestnut hair is a hair color which is a reddish shade of brown hair. In contrast to auburn hair, the reddish shade of chestnut is darker. Chestnut hair is common among the native peoples of Northern, Central, Western, and Eastern Europe and is also found in Asia Minor, West Asia and North Africa.
Auburn hair
Auburn hair ranges along a spectrum of light to dark red-brown shades. The chemicals which cause auburn hair are eumelanin (brown) and pheomelanin (red), with a higher proportion of red-causing pheomelanin than is found in average brown hair. It is most commonly found in individuals of Northern and Western European descent, but is extant in West and Central Asia and North Africa also. It can also be the result of a mutation in the melanocortin 1 receptor gene.
Red hair
Red hair ranges from light strawberry blond shades to titian, copper, and completely red. Red hair has the highest amounts of pheomelanin, around 67%, and usually low levels of eumelanin. At 1–2% of the west Eurasian population, it is the least common hair color in the world. It is most prominently found in the British Isles and in Udmurtia. Scotland has the highest proportion of redheads; 13 percent of the population has red hair and approximately 40 percent carry the recessive redhead gene. Red hair can also occur in Southern Europe, West Asia, North Africa and Central Asia.
Blond hair
Blond (sometimes blonde for women) hair ranges from pale white (platinum blond) to dark gold blond. Strawberry blond, a mixture of blond and red hair, is a much rarer type containing the most pheomelanin. Blond hair can have almost any proportion of pheomelanin and eumelanin, but has only small amounts of both. More pheomelanin creates a more golden or strawberry blond color, and more eumelanin creates an ash or sandy blond color. Blond hair is most commonly found in Northern and Northeastern Europeans and their descendants but can be found spread around most of Europe and also among West Asians and North Africans at lower frequencies. Studies in 2012 showed that naturally blond hair of Melanesians is caused by a recessive mutation in tyrosinase-related protein 1 (TYRP1). In the Solomon Islands, 26% of the population carry the gene; however, it is absent outside of Oceania.
Gray and white hair
Gray or white hair is not caused by a true gray or white pigment, but is due to a lack of pigmentation and melanin. The clear hairs appear as gray or white because of the way light is reflected from the hairs. Gray hair color typically occurs naturally as people age (see aging or achromotrichia below).
Marie Antoinette syndrome is a proposed phenomenon in which sudden whitening is caused by stress. It has been found that some hairs can become colored again when stress is reduced.
Conditions affecting hair color
Aging or achromotrichia
Children born with some hair colors may find it gradually darkens as they grow. Many blond, light brown, or red haired infants experience this. This is caused by genes being turned on and off during early childhood and puberty.
Changes in hair color typically occur naturally as people age, eventually turning the hair gray and then white. This is called achromotrichia. Achromotrichia normally begins in the early to mid-twenties in men and late twenties in women. More than 60 percent of Americans have some gray hair by age 40. The age at which graying begins seems almost entirely due to genetics. Sometimes people are born with gray hair because they inherit the trait.
The order in which graying happens is usually: nose hair, hair on the head, beard, body hair, eyebrows.
Hair coloring
Hair color can be changed by a chemical process. Hair coloring is classed as "permanent" or "semi-permanent".
Permanent hair color means that the hair's structure has been chemically altered until it is eventually cut away. This does not mean that the synthetic color will remain permanently. During the process, the natural color is removed, one or more shades, and synthetic color has been put in its place. All pigments wash out of the cuticle. Natural color stays in much longer and artificial will fade the fastest (depending on the color molecules and the form of the dye pigments).
Permanent hair color gives the most flexibility because it can make hair lighter or darker as well as changing tone and color, but there are negatives. Constant (monthly or six-weekly) maintenance is essential to match new hair growing in to the rest of the hair, and to remedy fading. A one-color permanent dye creates a flat, uniform color across the whole head, which can look unnatural and harsh, especially in a fair shade. To combat this, the modern trend is to use multiple colors—usually one color as a base with added highlights or lowlights in other shades.
Semi-permanent color washes out over a period of time—typically four to six weeks, so root regrowth is less noticeable. The final color of each strand is affected by its original color and porosity, so there will be subtle variations in color across the head—more natural and less harsh than a permanent dye. However, this means that gray and white hair will not dye to the same color as the rest of the head (in fact, some white hair will not absorb the color at all). A few gray and white hairs will blend in visually, but semi-permanent dye alone will not usually give the desired result where there is a lot of gray or white hair present. Sometimes a mixture of dyes is used while hair is greying: semi-permanent as a base color, with permanent highlights.
Semi-permanent hair color cannot lighten hair. Hair can only be lightened using chemical lighteners, such as bleach. Bleaching is always permanent because it removes the natural pigment.
"Rinses" are a form of temporary hair color, usually applied to hair during a shampoo and washed out again the next time the hair is washed.
| Biology and health sciences | Health and fitness: General | Health |
270925 | https://en.wikipedia.org/wiki/Bulletin%20board | Bulletin board | A bulletin board (pinboard, pin board, noticeboard, or notice board in British English) is a surface intended for the posting of public messages, for example, to advertise items wanted or for sale, announce events, or provide information. Bulletin boards are often made of a material such as cork to facilitate addition and removal of messages, as well as a writing surface such as blackboard or whiteboard. A bulletin board which combines a pinboard (corkboard) and writing surface is known as a combination bulletin board. Bulletin boards can also be entirely in the digital domain and placed on computer networks so people can leave and erase messages for other people to read and see, as in a bulletin board system.
Bulletin boards are particularly prevalent at universities. They are used by many sports groups and extracurricular groups and anything from local shops to official notices. Dormitory corridors, well-trafficked hallways, lobbies, and freestanding kiosks often have cork boards attached to facilitate the posting of notices. At some universities, lampposts, bollards, trees, and walls often become impromptu posting sites in areas where official boards are sparse in number.
Internet forums are a replacement for traditional bulletin boards. Online bulletin boards are sometimes referred to as message boards. The terms bulletin board, message board and even Internet forum are interchangeable, although often one bulletin board or message board can contain a number of Internet forums or discussion groups. An online board can serve the same purpose as a physical bulletin board, with the added benefit of not being bound by geographical location.
Magnet boards, or magnetic bulletin boards, are a popular substitute for cork boards because they lack the problem of board deterioration from the insertion and removal of pins over time.
History
1801: James Pillans, headmaster and geography teacher at the Old High School in Edinburgh, Scotland, is credited with inventing the first modern blackboard.
1925: George Brooks of Topeka, Kansas is issued a patent for the use of corkboard as a bulletin board which you could stick tacks into. The patent for George Brooks' invention, which would become a mainstay in homes and offices around the world, expired in 1941, which then allowed anyone to create and market their own versions of the product.
1940: George E. Fox, received a patent for a foam rubber pinboard with cardboard backing
1959: W.F. Lewis is issued a patent for a combined chalkboard and bulletin board.
1978: The concept of the bulletin board entered the information age when software developers Ward Christensen and Randy Suess launched the first public dialup bulletin board system.
2010: Digital signage displays started to replace bulletin boards as a means to reduce clutter and provide real-time information.
2010: Pinterest, a content and photo-sharing website meant to serve as online personal bulletin boards, is founded by Ben Silbermann, Paul Sciarra and Evan Sharp.
2017: General Enchantment is issued a utility patent for a physical electronic bulletin board system, that includes a physical writing or pinning surface and an electronic display, like a mobile computing device (tablet computer), capable of running (digital signage) software that augments the sharing of analogue information with digital content.
| Technology | Media and communication: Basics | null |
271046 | https://en.wikipedia.org/wiki/Resonance%20%28chemistry%29 | Resonance (chemistry) | In chemistry, resonance, also called mesomerism, is a way of describing bonding in certain molecules or polyatomic ions by the combination of several contributing structures (or forms, also variously known as resonance structures or canonical structures) into a resonance hybrid (or hybrid structure) in valence bond theory. It has particular value for analyzing delocalized electrons where the bonding cannot be expressed by one single Lewis structure. The resonance hybrid is the accurate structure for a molecule or ion; it is an average of the theoretical (or hypothetical) contributing structures.
Overview
Under the framework of valence bond theory, resonance is an extension of the idea that the bonding in a chemical species can be described by a Lewis structure. For many chemical species, a single Lewis structure, consisting of atoms obeying the octet rule, possibly bearing formal charges, and connected by bonds of positive integer order, is sufficient for describing the chemical bonding and rationalizing experimentally determined molecular properties like bond lengths, angles, and dipole moment. However, in some cases, more than one Lewis structure could be drawn, and experimental properties are inconsistent with any one structure. In order to address this type of situation, several contributing structures are considered together as an average, and the molecule is said to be represented by a resonance hybrid in which several Lewis structures are used collectively to describe its true structure. For instance, in NO2–, nitrite anion, the two N–O bond lengths are equal, even though no single Lewis structure has two N–O bonds with the same formal bond order. However, its measured structure is consistent with a description as a resonance hybrid of the two major contributing structures shown above: it has two equal N–O bonds of 125 pm, intermediate in length between a typical N–O single bond (145 pm in hydroxylamine, H2N–OH) and N–O double bond (115 pm in nitronium ion, [O=N=O]+). According to the contributing structures, each N–O bond is an average of a formal single and formal double bond, leading to a true bond order of 1.5. By virtue of this averaging, the Lewis description of the bonding in NO2– is reconciled with the experimental fact that the anion has equivalent N–O bonds.
The resonance hybrid represents the actual molecule as the "average" of the contributing structures, with bond lengths and partial charges taking on intermediate values compared to those expected for the individual Lewis structures of the contributors, were they to exist as "real" chemical entities. The contributing structures differ only in the formal apportionment of electrons to the atoms, and not in the actual physically and chemically significant electron or spin density. While contributing structures may differ in formal bond orders and in formal charge assignments, all contributing structures must have the same number of valence electrons and the same spin multiplicity.
Because electron delocalization lowers the potential energy of a system, any species represented by a resonance hybrid is more stable than any of the (hypothetical) contributing structures. Electron delocalization stabilizes a molecule because the electrons are more evenly spread out over the molecule, decreasing electron-electron repulsion. The difference in potential energy between the actual species and the (computed) energy of the contributing structure with the lowest potential energy is called the resonance energy or delocalization energy. The magnitude of the resonance energy depends on assumptions made about the hypothetical "non-stabilized" species and the computational methods used and does not represent a measurable physical quantity, although comparisons of resonance energies computed under similar assumptions and conditions may be chemically meaningful.
Molecules with an extended π system such as linear polyenes and polyaromatic compounds are well described by resonance hybrids as well as by delocalised orbitals in molecular orbital theory.
Resonance vs isomerism
Resonance is to be distinguished from isomerism. Isomers are molecules with the same chemical formula but are distinct chemical species with different arrangements of atomic nuclei in space. Resonance contributors of a molecule, on the other hand, can only differ in the way electrons are formally assigned to atoms in the Lewis structure depictions of the molecule. Specifically, when a molecular structure is said to be represented by a resonance hybrid, it does not mean that electrons of the molecule are "resonating" or shifting back and forth between several sets of positions, each one represented by a Lewis structure. Rather, it means that the set of contributing structures represents an intermediate structure (a weighted average of the contributors), with a single, well-defined geometry and distribution of electrons. It is incorrect to regard resonance hybrids as rapidly interconverting isomers, even though the term "resonance" might evoke such an image. (As described below, the term "resonance" originated as a classical physics analogy for a quantum mechanical phenomenon, so it should not be construed too literally.) Symbolically, the double headed arrow A<->B is used to indicate that A and B are contributing forms of a single chemical species (as opposed to an equilibrium arrow, e.g., A <=> B; see below for details on usage).
A non-chemical analogy is illustrative: one can describe the characteristics of a real animal, the narwhal, in terms of the characteristics of two mythical creatures: the unicorn, a creature with a single horn on its head, and the leviathan, a large, whale-like creature. The narwhal is not a creature that goes back and forth between being a unicorn and being a leviathan, nor do the unicorn and leviathan have any physical existence outside the collective human imagination. Nevertheless, describing the narwhal in terms of these imaginary creatures provides a reasonably good description of its physical characteristics.
Due to confusion with the physical meaning of the word resonance, as no entities actually physically "resonate", it has been suggested that the term resonance be abandoned in favor of delocalization and resonance energy abandoned in favor of delocalization energy. A resonance structure becomes a contributing structure and the resonance hybrid becomes the hybrid structure. The double headed arrows would be replaced by commas to illustrate a set of structures, as arrows of any type may suggest that a chemical change is taking place.
Representation in diagrams
In diagrams, contributing structures are typically separated by double-headed arrows (↔). The arrow should not be confused with the right and left pointing equilibrium arrow (). All structures together may be enclosed in large square brackets, to indicate they picture one single molecule or ion, not different species in a chemical equilibrium.
Alternatively to the use of contributing structures in diagrams, a hybrid structure can be used. In a hybrid structure, pi bonds that are involved in resonance are usually pictured as curves or dashed lines, indicating that these are partial rather than normal complete pi bonds. In benzene and other aromatic rings, the delocalized pi-electrons are sometimes pictured as a solid circle.
History
The concept first appeared in 1899 in Johannes Thiele's "Partial Valence Hypothesis" to explain the unusual stability of benzene which would not be expected from August Kekulé's structure proposed in 1865 with alternating single and double bonds. Benzene undergoes substitution reactions, rather than addition reactions as typical for alkenes. He proposed that the carbon-carbon bond in benzene is intermediate of a single and double bond.
The resonance proposal also helped explain the number of isomers of benzene derivatives. For example, Kekulé's structure would predict four dibromobenzene isomers, including two ortho isomers with the brominated carbon atoms joined by either a single or a double bond. In reality there are only three dibromobenzene isomers and only one is ortho, in agreement with the idea that there is only one type of carbon-carbon bond, intermediate between a single and a double bond.
The mechanism of resonance was introduced into quantum mechanics by Werner Heisenberg in 1926 in a discussion of the quantum states of the helium atom. He compared the structure of the helium atom with the classical system of resonating coupled harmonic oscillators. In the classical system, the coupling produces two modes, one of which is lower in frequency than either of the uncoupled vibrations; quantum mechanically, this lower frequency is interpreted as a lower energy. Linus Pauling used this mechanism to explain the partial valence of molecules in 1928, and developed it further in a series of papers in 1931-1933. The alternative term mesomerism popular in German and French publications with the same meaning was introduced by C. K. Ingold in 1938, but did not catch on in the English literature. The current concept of mesomeric effect has taken on a related but different meaning. The double headed arrow was introduced by the German chemist Fritz Arndt who preferred the German phrase zwischenstufe or intermediate stage.
Resonance theory dominated over competing Hückel method for two decades thanks to being relatively easier to understand for chemists without fundamental physics background, even if they couldn't grasp the concept of quantum superposition and confused it with tautomerism. Pauling and Wheland themselves characterized Erich Hückel's approach as "cumbersome" at the time, and his lack of communication skills contributed: when Robert Robinson sent him a friendly request, he responded arrogantly that he is not interested in organic chemistry.
In the Soviet Union, resonance theory – especially as developed by Pauling – was attacked in the early 1950s as being contrary to the Marxist principles of dialectical materialism, and in June 1951 the Soviet Academy of Sciences under the leadership of Alexander Nesmeyanov convened a conference on the chemical structure of organic compounds, attended by 400 physicists, chemists, and philosophers, where "the pseudo-scientific essence of the theory of resonance was exposed and unmasked".
Major and minor contributors
One contributing structure may resemble the actual molecule more than another (in the sense of energy and stability). Structures with a low value of potential energy are more stable than those with high values and resemble the actual structure more. The most stable contributing structures are called major contributors. Energetically unfavourable and therefore less favorable structures are minor contributors. With rules listed in rough order of diminishing importance, major contributors are generally structures that
obey as much as possible the octet rule (8 valence electrons around each atom rather than having deficiencies or surplus, or 2 electrons for Period 1 elements);
have a maximum number of covalent bonds;
carry a minimum of formally charged atoms, with the separation for unlike and like charges minimized and maximized, respectively;
place negative charge, if any, on the most electronegative atoms and positive charge, if any, on the most electropositive;
do not deviate substantially from idealized bond lengths and angles (e.g., the relative unimportance of Dewar-type resonance contributors for benzene);
maintain aromatic substructures locally while avoiding anti-aromatic ones (see Clar sextet and biphenylene).
A maximum of eight valence electrons is strict for the Period 2 elements Be, B, C, N, O, and F, as is a maximum of two for H and He and effectively for Li as well. The issue of expansion of the valence shell of third period and heavier main group elements is controversial. A Lewis structure in which a central atom has a valence electron count greater than eight traditionally implies the participation of d orbitals in bonding. However, the consensus opinion is that while they may make a marginal contribution, the participation of d orbitals is unimportant, and the bonding of so-called hypervalent molecules are, for the most part, better explained by charge-separated contributing forms that depict three-center four-electron bonding. Nevertheless, by tradition, expanded octet structures are still commonly drawn for functional groups like sulfoxides, sulfones, and phosphorus ylides, for example. Regarded as a formalism that does not necessarily reflect the true electronic structure, such depictions are preferred by the IUPAC over structures featuring partial bonds, charge separation, or dative bonds.
Equivalent contributors contribute equally to the actual structure, while the importance of nonequivalent contributors is determined by the extent to which they conform to the properties listed above. A larger number of significant contributing structures and a more voluminous space available for delocalized electrons lead to stabilization (lowering of the energy) of the molecule.
Examples
Aromatic molecules
In benzene the two cyclohexatriene Kekulé structures, first proposed by Kekulé, are taken together as contributing structures to represent the total structure. In the hybrid structure on the right, the dashed hexagon replaces three double bonds, and represents six electrons in a set of three molecular orbitals of π symmetry, with a nodal plane in the plane of the molecule.
In furan a lone pair of the oxygen atom interacts with the π orbitals of the carbon atoms. The curved arrows depict the permutation of delocalized π electrons, which results in different contributors.
Electron-rich molecules
The ozone molecule is represented by two contributing structures. In reality the two terminal oxygen atoms are equivalent and the hybrid structure is drawn on the right with a charge of − on both oxygen atoms and partial double bonds with a full and dashed line and bond order .
For hypervalent molecules, the rationalization described above can be applied to generate contributing structures to explain the bonding in such molecules. Shown below are the contributing structures of a 3c-4e bond in xenon difluoride.
[\mathsf{F-XeF^- <-> F^-Xe-F}]
Electron-deficient molecules
The allyl cation has two contributing structures with a positive charge on the terminal carbon atoms. In the hybrid structure their charge is +. The full positive charge can also be depicted as delocalized among three carbon atoms.
The diborane molecule is described by contributing structures, each with electron-deficiency on different atoms. This reduces the electron-deficiency on each atom and stabilizes the molecule. Below are the contributing structures of an individual 3c-2e bond in diborane.
Reactive intermediates
Often, reactive intermediates such as carbocations and free radicals have more delocalized structure than their parent reactants, giving rise to unexpected products. The classical example is allylic rearrangement. When 1 mole of HCl adds to 1 mole of 1,3-butadiene, in addition to the ordinarily expected product 3-chloro-1-butene, we also find 1-chloro-2-butene. Isotope labelling experiments have shown that what happens here is that the additional double bond shifts from 1,2 position to 2,3 position in some of the product. This and other evidence (such as NMR in superacid solutions) shows that the intermediate carbocation must have a highly delocalized structure, different from its mostly classical (delocalization exists but is small) parent molecule. This cation (an allylic cation) can be represented using resonance, as shown above.
This observation of greater delocalization in less stable molecules is quite general. The excited states of conjugated dienes are stabilised more by conjugation than their ground states, causing them to become organic dyes.
A well-studied example of delocalization that does not involve π electrons (hyperconjugation) can be observed in the non-classical 2-Norbornyl cation Another example is methanium (). These can be viewed as containing three-center two-electron bonds and are represented either by contributing structures involving rearrangement of σ electrons or by a special notation, a Y that has the three nuclei at its three points.
Delocalized electrons are important for several reasons; a major one is that an expected chemical reaction may not occur because the electrons delocalize to a more stable configuration, resulting in a reaction that happens at a different location. An example is the Friedel–Crafts alkylation of benzene with 1-chloro-2-methylpropane; the carbocation rearranges to a tert-butyl group stabilized by hyperconjugation, a particular form of delocalization.
Benzene
Bond lengths
Comparing the two contributing structures of benzene, all single and double bonds are interchanged. Bond lengths can be measured, for example using X-ray diffraction. The average length of a C–C single bond is 154 pm; that of a C=C double bond is 133 pm. In localized cyclohexatriene, the carbon–carbon bonds should be alternating 154 and 133 pm. Instead, all carbon–carbon bonds in benzene are found to be about 139 pm, a bond length intermediate between single and double bond. This mixed single and double bond (or triple bond) character is typical for all molecules in which bonds have a different bond order in different contributing structures. Bond lengths can be compared using bond orders. For example, in cyclohexane the bond order is 1 while that in benzene is 1 + (3 ÷ 6) = . Consequently, benzene has more double bond character and hence has a shorter bond length than cyclohexane.
Resonance energy
Resonance (or delocalization) energy is the amount of energy needed to convert the true delocalized structure into that of the most stable contributing structure. The empirical resonance energy can be estimated by comparing the enthalpy change of hydrogenation of the real substance with that estimated for the contributing structure.
The complete hydrogenation of benzene to cyclohexane via 1,3-cyclohexadiene and cyclohexene is exothermic; 1 mole of benzene delivers 208.4 kJ (49.8 kcal).
Hydrogenation of one mole of double bonds delivers 119.7 kJ (28.6 kcal), as can be deduced from the last step, the hydrogenation of cyclohexene. In benzene, however, 23.4 kJ (5.6 kcal) are needed to hydrogenate one mole of double bonds. The difference, being 143.1 kJ (34.2 kcal), is the empirical resonance energy of benzene. Because 1,3-cyclohexadiene also has a small delocalization energy (7.6 kJ or 1.8 kcal/mol) the net resonance energy, relative to the localized cyclohexatriene, is a bit higher: 151 kJ or 36 kcal/mol.
This measured resonance energy is also the difference between the hydrogenation energy of three 'non-resonance' double bonds and the measured hydrogenation energy:
(3 × 119.7) − 208.4 = 150.7 kJ/mol (36 kcal).
Regardless of their exact values, resonance energies of various related compounds provide insights into their bonding. The resonance energies for pyrrole, thiophene, and furan are, respectively, 88, 121, and
67 kJ/mol (21, 29, and 16 kcal/mol). Thus, these heterocycles are far less aromatic than benzene, as is manifested in the lability of these rings.
Quantum mechanical description in valence bond (VB) theory
Resonance has a deeper significance in the mathematical formalism of valence bond theory (VB). Quantum mechanics requires that the wavefunction of a molecule obey its observed symmetry. If a single contributing structure does not achieve this, resonance is invoked.
For example, in benzene, valence bond theory begins with the two Kekulé structures which do not individually possess the sixfold symmetry of the real molecule. The theory constructs the actual wave function as a linear superposition of the wave functions representing the two structures. As both Kekulé structures have equal energy, they are equal contributors to the overall structure – the superposition is an equally weighted average, or a 1:1 linear combination of the two in the case of benzene. The symmetric combination gives the ground state, while the antisymmetric combination gives the first excited state, as shown.
In general, the superposition is written with undetermined coefficients, which are then variationally optimized to find the lowest possible energy for the given set of basis wave functions. When more contributing structures are included, the molecular wave function becomes more accurate and more excited states can be derived from different combinations of the contributing structures.
Comparison with molecular orbital (MO) theory
In molecular orbital theory, the main alternative to valence bond theory, the molecular orbitals (MOs) are approximated as sums of all the atomic orbitals (AOs) on all the atoms; there are as many MOs as AOs. Each AOi has a weighting coefficient ci that indicates the AO's contribution to a particular MO. For example, in benzene, the MO model gives us 6 π MOs which are combinations of the 2pz AOs on each of the 6 C atoms. Thus, each π MO is delocalized over the whole benzene molecule and any electron occupying an MO will be delocalized over the whole molecule. This MO interpretation has inspired the picture of the benzene ring as a hexagon with a circle inside. When describing benzene, the VB concept of localized σ bonds and the MO concept of delocalized π orbitals are frequently combined in elementary chemistry courses.
The contributing structures in the VB model are particularly useful in predicting the effect of substituents on π systems such as benzene. They lead to the models of contributing structures for an electron-withdrawing group and electron-releasing group on benzene. The utility of MO theory is that a quantitative indication of the charge from the π system on an atom can be obtained from the squares of the weighting coefficient ci on atom Ci. Charge qi ≈ c. The reason for squaring the coefficient is that if an electron is described by an AO, then the square of the AO gives the electron density. The AOs are adjusted (normalized) so that AO2 = 1, and qi ≈ (ciAOi)2 ≈ c. In benzene, qi = 1 on each C atom. With an electron-withdrawing group qi < 1 on the ortho and para C atoms and qi > 1 for an electron-releasing group.
Coefficients
Weighting of the contributing structures in terms of their contribution to the overall structure can be calculated in multiple ways, using "Ab initio" methods derived from Valence Bond theory, or else from the Natural Bond Orbitals (NBO) approaches of Weinhold NBO5 , or finally from empirical calculations based on the Hückel method. A Hückel method-based software for teaching resonance is available on the HuLiS Web site.
Charge delocalization
In the case of ions it is common to speak about delocalized charge (charge delocalization). An example of delocalized charge in ions can be found in the carboxylate group, wherein the negative charge is centered equally on the two oxygen atoms. Charge delocalization in anions is an important factor determining their reactivity (generally: the higher the extent of delocalization the lower the reactivity) and, specifically, the acid strength of their conjugate acids. As a general rule, the better delocalized is the charge in an anion the stronger is its conjugate acid. For example, the negative charge in perchlorate anion () is evenly distributed among the symmetrically oriented oxygen atoms (and a part of it is also kept by the central chlorine atom). This excellent charge delocalization combined with the high number of oxygen atoms (four) and high electronegativity of the central chlorine atom leads to perchloric acid being one of the strongest known acids with a pKa value of −10.
The extent of charge delocalization in an anion can be quantitatively expressed via the WAPS (weighted average positive sigma) parameter parameter and an analogous WANS (weighted average negative sigma) parameter is used for cations.
WAPS and WANS values are given in e/Å4. Larger values indicate more localized charge in the corresponding ion.
| Physical sciences | Bond structure | Chemistry |
271396 | https://en.wikipedia.org/wiki/Pennyweight | Pennyweight | A pennyweight (dwt) is a unit of mass equal to 24 grains, of a troy ounce, of a troy pound, approximately 0.054857 avoirdupois ounce and exactly 1.55517384 grams. It is abbreviated dwt, d standing for denarius (an ancient Roman coin), and later used as the symbol of an old British penny (see £sd).
History
In the Middle Ages, an English penny's weight was literally, as well as monetarily, of an ounce and of a pound of sterling silver. At that time, the pound unit in use in England was the Tower pound, equal to 7,680 Tower grains (also known as wheat grains). The medieval English pennyweight was thus equal to 32 Tower grains. When Troy weights replaced Tower weights in 1527, the Troy weights were defined in such a way that the old Tower pound came out to exactly 5,400 Troy grains (also known as barleycorns), the Tower pennyweight Troy grains (and thus approximately 1.46 grams). After 1527, the English pennyweight was the Troy pennyweight. of 24 Troy grains. Thus the Troy pound, ounce, and pennyweight, with their definitions given in terms of the Troy grain instead of in terms of the Tower grain, were or 6.667% more than the Tower equivalents.
Usage
The troy pound and the pennyweight lost their official status in the United Kingdom in the Weights and Measures Act of 1878; only the troy ounce and its decimal subdivisions remained official. The troy ounce enjoys a specific legal exemption from metrication in the UK.
The pennyweight is the common weight used in the valuation and measurement of precious metals. Jewellers use the pennyweight in calculating the amount and cost of precious metals used in fabricating or casting jewellery. Similarly, dentists and dental labs still use the pennyweight as the measure of precious metals in dental crowns and inlays.
Pennyweight and grains are still used to weigh gooseberries in competitions in Cheshire, northwest UK. Over the Pennines in Yorkshire the alternative drams and grains measurement has been used since a new set of scales was purchased by the Egton Bridge Old Gooseberry Society in 1937. As of 2018, the world record for the heaviest gooseberry of () was held by Kelvin Archer of Cheshire.
The most common abbreviation for pennyweight is dwt; d, for the Roman denarius, was the abbreviation for penny before Decimalisation of the British monetary system. Alternate abbreviations are pwt and PW.
Uses unrelated to weight
Although the abbreviations are the same, the pennyweight bears no relation to the weight of the American penny nail. That name is derived from the price for a hundred nails in 15th century England: the larger the nail, the higher the cost per hundred.
The pennyweight also bears no relation to the weight of the American "penny" (1 cent) coin, which weighs 2.5 g (for those minted after 1982).
Conversion
| Physical sciences | Mass and weight | Basics and measurement |
3572234 | https://en.wikipedia.org/wiki/Terrace%20%28geology%29 | Terrace (geology) | In geology, a terrace is a step-like landform. A terrace consists of a flat or gently sloping geomorphic surface, called a tread, that is typically bounded on one side by a steeper ascending slope, which is called a "riser" or "scarp". The tread and the steeper descending slope (riser or scarp) together constitute the terrace. Terraces can also consist of a tread bounded on all sides by a descending riser or scarp. A narrow terrace is often called a bench.
The sediments underlying the tread and riser of a terrace are also commonly, but incorrectly, called terraces, leading to confusion.
Terraces are formed in various ways.
Fluvial terraces
Fluvial terraces are remnants of the former floodplain of a stream or river. They are formed by the downcutting of a river or stream channel into and the abandonment and lateral erosion of its former floodplain. The downcutting, abandonment, and lateral erosion of a former floodplain can be the result of either changes in sea level, local or regional tectonic uplift; changes in local or regional climate; changes in the amount of sediment being carried by the river or stream; change in discharge of the river; or a complex mixture of these and other factors. The most common sources of the variations in rivers and streams that create fluvial terraces are vegetative, geomorphic, and hydrologic responses to climate. More recently, the direct modification of rivers and streams and their watersheds by cultural processes have resulted in the development of terraces along many rivers and streams.
Kame terraces
Kame terraces are formed on the side of a glacial valley and are the deposits of meltwater streams flowing between the ice and the adjacent valley side.
Marine terraces
A marine terrace represents the former shoreline of a sea or ocean. It can be formed by marine abrasion or erosion of materials comprising the shoreline (marine-cut terraces or wave-cut platforms); the accumulations of sediments in the shallow-water to slightly emerged coastal environments (marine-built terraces or raised beach); or the bioconstruction by coral reefs and accumulation of reef materials (reef flats) in intertropical regions.
The formation of a marine terrace follows this general process: A wave cut platform must be carved into bedrock (high wave energy is needed for this process). Although this is the first step to the process for the formation of a marine terrace, not all wave cut platforms will become a marine terrace. After the wave cut platform is formed it must be removed from interaction with the high wave energy. This process happens by either change in sea level due to glacial-interglacial cycles or tectonically rising landmasses. When the wave cut has been raised above sea level it is preserved. The terraces are most commonly preserved in flights along the coastline.
Lacustrine terraces
A lake (lacustrine) terrace represents the former shoreline of either a nonglacial, glacial, or proglacial lake. As with marine terraces, a lake terrace can be formed by either the abrasion or erosion of materials comprising the shoreline, the accumulations of sediments in the shallow-water to slightly emerged environments, or some combination of these. Given the smaller size of lakes relative to the size of typical marine water bodies, lake terraces are overall significantly narrower and less well developed than marine terraces. However, not all lake terraces are relict shorelines. In case of the lake terraces of ancient ice-walled lakes, some proglacial lakes, and alluvium-dammed (slackwater) lakes, they often represent the relict bottom of these lakes. Finally, glaciolacustrine kame terraces are either the relict deltas or bottoms of ancient ice marginal lakes.
Structural terraces
In geomorphology, a structural terrace is a terrace created by the differential erosion of flat-lying or nearly flat-lying layered strata. The terrace results from preferential stripping by erosion of a layer of softer strata from an underlying layer of harder strata. The preferential removal of softer material exposes the flat surface of the underlying harder layer, creating the tread of a structural terrace. Structural terraces are commonly paired and not always associated with river valleys.
Travertine terraces
A travertine terrace is formed when geothermally heated supersaturated alkaline waters emerge to the surface and form waterfalls of precipitated carbonates.
| Physical sciences | Other erosional landforms | Earth science |
3573584 | https://en.wikipedia.org/wiki/Modelling%20clay | Modelling clay | Modelling clay or modelling compound is any of a group of malleable substances used in building and sculpting. The material compositions and production processes vary considerably.
Ceramic clay
Ceramic clays are water-based substances made from clay minerals and other raw materials. They are baked at high temperatures in a process known as firing to create ceramics, such as terra cotta, earthenware, stoneware, and porcelain. Paper clay produced by pottery clay manufacturers is a clay body to which a small percentage of processed cellulose fiber has been added. When kiln-fired, the paper burns out, leaving the clay body. Consequently, the firing temperatures and glazes selection should be the same on those used with the clay body.
Oil-based clay
Oil-based clays are made from combinations of oils, waxes, and clay minerals.
Unlike water, the oils do not evaporate and oil-based clays remain malleable even when left in dry environments for long periods. Articles made from oil-based clays cannot be fired, and thus are not ceramics. Because rising temperature decreases oil viscosity, the malleability is influenced by heating or cooling the clay. Oil-based clay is not water-soluble. As it can be re-used, it is a material commonly used by stop motion animators who need to bend and move their models. It is available in a multitude of colors and is non-toxic. Readily worked in fine detail, oil-based clays are also suitable for the creation of detailed sculptures from which molds can be made. Castings and reproductions can then be produced from much more durable materials. Cars and airplanes may be created using industrial design-grade modelling clay.
Oil-based clays are referred to by multiple of genericized trademarks.
Plastilin (or Plasteline), which was patented in Germany by Franz Kolb in 1880, was developed by Claude Chavant in 1892 and trademarked in 1927.
Plasticine was invented in 1897 by William Harbutt of Bathampton, England.
Plastilina is trademarked as Roma Plastilina by Sculpture House, Inc. According to their website, their formula is 100 years old. Roma Plastilina contains sulfur, and since certain moldmaking compounds do not set in sulfur's presence, making molds of items made of industrial plasticine is difficult.
Polymer clay
Polymer clay is a modelling material that cures when heated from for 15 minutes per of thickness, and does not significantly shrink or change shape during the process. Despite being called "clay", it generally contains no clay minerals. Polymer clay is sold in craft, hobby, and art stores, and is used by artists, hobbyists, and children. Polymer clay is used in animation, since it allows static forms to be manipulated frame after frame. Leading brands of polymer clay include Fimo, Kato Polyclay, Sculpey, Modello and Crafty Argentina.
Paper clay
Paper clay is handmade or commercially available clay to which a small percentage of processed cellulose fiber is added. The fiber increases the tensile strength of the dry clay and enables dry-to-dry and wet-to-dry joins. Commercial paper clays air-dry to a firm, lightweight sculpture, with minimal shrinking during the drying process.
Paper clay can be used as an unfired body in craft and doll-making. It is used in ceramic art studios as sculptural and functional studio pottery.
| Technology | Artist's tools | null |
8282374 | https://en.wikipedia.org/wiki/Tropical%20cyclone | Tropical cyclone | A tropical cyclone is a rapidly rotating storm system with a low-pressure area, a closed low-level atmospheric circulation, strong winds, and a spiral arrangement of thunderstorms that produce heavy rain and squalls. Depending on its location and strength, a tropical cyclone is called a hurricane (), typhoon (), tropical storm, cyclonic storm, tropical depression, or simply cyclone. A hurricane is a strong tropical cyclone that occurs in the Atlantic Ocean or northeastern Pacific Ocean. A typhoon occurs in the northwestern Pacific Ocean. In the Indian Ocean and South Pacific, comparable storms are referred to as "tropical cyclones". In modern times, on average around 80 to 90 named tropical cyclones form each year around the world, over half of which develop hurricane-force winds of or more.
Tropical cyclones typically form over large bodies of relatively warm water. They derive their energy through the evaporation of water from the ocean surface, which ultimately condenses into clouds and rain when moist air rises and cools to saturation. This energy source differs from that of mid-latitude cyclonic storms, such as nor'easters and European windstorms, which are powered primarily by horizontal temperature contrasts. Tropical cyclones are typically between in diameter. The strong rotating winds of a tropical cyclone are a result of the conservation of angular momentum imparted by the Earth's rotation as air flows inwards toward the axis of rotation. As a result, cyclones rarely form within 5° of the equator. Tropical cyclones are very rare in the South Atlantic (although occasional examples do occur) due to consistently strong wind shear and a weak Intertropical Convergence Zone. In contrast, the African easterly jet and areas of atmospheric instability give rise to cyclones in the Atlantic Ocean and Caribbean Sea.
Heat energy from the ocean acts as the accelerator for tropical cyclones. This causes inland regions to suffer far less damage from cyclones than coastal regions, although the impacts of flooding are felt across the board. Coastal damage may be caused by strong winds and rain, high waves (due to winds), storm surges (due to wind and severe pressure changes), and the potential of spawning tornadoes. Climate change affects tropical cyclones in several ways. Scientists found that climate change can exacerbate the impact of tropical cyclones by increasing their duration, occurrence, and intensity due to the warming of ocean waters and intensification of the water cycle. Tropical cyclones draw in air from a large area and concentrate the water content of that air into precipitation over a much smaller area. This replenishing of moisture-bearing air after rain may cause multi-hour or multi-day extremely heavy rain up to from the coastline, far beyond the amount of water that the local atmosphere holds at any one time. This in turn can lead to river flooding, overland flooding, and a general overwhelming of local water control structures across a large area.
Definition and terminology
A tropical cyclone is the generic term for a warm-cored, non-frontal synoptic-scale low-pressure system over tropical or subtropical waters around the world. The systems generally have a well-defined center which is surrounded by deep atmospheric convection and a closed wind circulation at the surface. A tropical cyclone is generally deemed to have formed once mean surface winds in excess of are observed. It is assumed at this stage that a tropical cyclone has become self-sustaining and can continue to intensify without any help from its environment.
Depending on its location and strength, a tropical cyclone is referred to by different names, including hurricane, typhoon, tropical storm, cyclonic storm, tropical depression, or simply cyclone. A hurricane is a strong tropical cyclone that occurs in the Atlantic Ocean or northeastern Pacific Ocean, and a typhoon occurs in the northwestern Pacific Ocean. In the Indian Ocean and South Pacific, comparable storms are referred to as "tropical cyclones", and such storms in the Indian Ocean can also be called "severe cyclonic storms".
Tropical refers to the geographical origin of these systems, which form almost exclusively over tropical seas. Cyclone refers to their winds moving in a circle, whirling round their central clear eye, with their surface winds blowing counterclockwise in the Northern Hemisphere and clockwise in the Southern Hemisphere. The opposite direction of circulation is due to the Coriolis effect.
Formation
Tropical cyclones tend to develop during the summer, but have been noted in nearly every month in most tropical cyclone basins. Tropical cyclones on either side of the Equator generally have their origins in the Intertropical Convergence Zone, where winds blow from either the northeast or southeast. Within this broad area of low-pressure, air is heated over the warm tropical ocean and rises in discrete parcels, which causes thundery showers to form. These showers dissipate quite quickly; however, they can group together into large clusters of thunderstorms. This creates a flow of warm, moist, rapidly rising air, which starts to rotate cyclonically as it interacts with the rotation of the earth.
Several factors are required for these thunderstorms to develop further, including sea surface temperatures of around and low vertical wind shear surrounding the system, atmospheric instability, high humidity in the lower to middle levels of the troposphere, enough Coriolis force to develop a low-pressure center, and a pre-existing low-level focus or disturbance.
There is a limit on tropical cyclone intensity which is strongly related to the water temperatures along its path. and upper-level divergence.
An average of 86 tropical cyclones of tropical storm intensity form annually worldwide. Of those, 47 reach strength higher than , and 20 become intense tropical cyclones, of at least Category 3 intensity on the Saffir–Simpson scale.
Climate oscillations such as El Niño–Southern Oscillation (ENSO) and the Madden–Julian oscillation modulate the timing and frequency of tropical cyclone development. Rossby waves can aid in the formation of a new tropical cyclone by disseminating the energy of an existing, mature storm. Kelvin waves can contribute to tropical cyclone formation by regulating the development of the westerlies. Cyclone formation is usually reduced 3 days prior to the wave's crest and increased during the 3 days after.
Formation regions and warning centers
The majority of tropical cyclones each year form in one of seven tropical cyclone basins, which are monitored by a variety of meteorological services and warning centers. Ten of these warning centers worldwide are designated as either a Regional Specialized Meteorological Centre or a Tropical Cyclone Warning Centre by the World Meteorological Organization's (WMO) tropical cyclone programme. These warning centers issue advisories which provide basic information and cover a systems present, forecast position, movement and intensity, in their designated areas of responsibility.
Meteorological services around the world are generally responsible for issuing warnings for their own country. There are exceptions, as the United States National Hurricane Center and Fiji Meteorological Service issue alerts, watches and warnings for various island nations in their areas of responsibility. The United States Joint Typhoon Warning Center and Fleet Weather Center also publicly issue warnings about tropical cyclones on behalf of the United States Government. The Brazilian Navy Hydrographic Center names South Atlantic tropical cyclones, however the South Atlantic is not a major basin, and not an official basin according to the WMO.
Interactions with climate
Each year on average, around 80 to 90 named tropical cyclones form around the world, of which over half develop hurricane-force winds of or more. Worldwide, tropical cyclone activity peaks in late summer, when the difference between temperatures aloft and sea surface temperatures is the greatest. However, each particular basin has its own seasonal patterns. On a worldwide scale, May is the least active month, while September is the most active month. November is the only month in which all the tropical cyclone basins are in season.
In the Northern Atlantic Ocean, a distinct cyclone season occurs from June 1 to November 30, sharply peaking from late August through September. The statistical peak of the Atlantic hurricane season is September 10.
The Northeast Pacific Ocean has a broader period of activity, but in a similar time frame to the Atlantic. The Northwest Pacific sees tropical cyclones year-round, with a minimum in February and March and a peak in early September. In the North Indian basin, storms are most common from April to December, with peaks in May and November. In the Southern Hemisphere, the tropical cyclone year begins on July 1 and runs all year-round encompassing the tropical cyclone seasons, which run from November 1 until the end of April, with peaks in mid-February to early March.
Of various modes of variability in the climate system, El Niño–Southern Oscillation has the largest effect on tropical cyclone activity. Most tropical cyclones form on the side of the subtropical ridge closer to the equator, then move poleward past the ridge axis before recurving into the main belt of the Westerlies. When the subtropical ridge position shifts due to El Niño, so will the preferred tropical cyclone tracks. Areas west of Japan and Korea tend to experience much fewer September–November tropical cyclone impacts during El Niño and neutral years.
During La Niña years, the formation of tropical cyclones, along with the subtropical ridge position, shifts westward across the western Pacific Ocean, which increases the landfall threat to China and much greater intensity in the Philippines. The Atlantic Ocean experiences depressed activity due to increased vertical wind shear across the region during El Niño years. Tropical cyclones are further influenced by the Atlantic Meridional Mode, the Quasi-biennial oscillation and the Madden–Julian oscillation.
Influence of climate change
The IPCC Sixth Assessment Report summarize the latest scientific findings about the impact of climate change on tropical cyclones. According to the report, we have now better understanding about the impact of climate change on tropical storm than before. Major tropical storms likely became more frequent in the last 40 years. We can say with high confidence that climate change increase rainfall during tropical cyclones. We can say with high confidence that a 1.5 degree warming lead to "increased proportion of and peak wind speeds of intense tropical cyclones". We can say with medium confidence that regional impacts of further warming include more intense tropical cyclones and/or extratropical storms.
Climate change can affect tropical cyclones in a variety of ways: an intensification of rainfall and wind speed, a decrease in overall frequency, an increase in the frequency of very intense storms and a poleward extension of where the cyclones reach maximum intensity are among the possible consequences of human-induced climate change. Tropical cyclones use warm, moist air as their fuel. As climate change is warming ocean temperatures, there is potentially more of this fuel available.
Between 1979 and 2017, there was a global increase in the proportion of tropical cyclones of Category 3 and higher on the Saffir–Simpson scale. The trend was most clear in the North Atlantic and in the Southern Indian Ocean. In the North Pacific, tropical cyclones have been moving poleward into colder waters and there was no increase in intensity over this period. With warming, a greater percentage (+13%) of tropical cyclones are expected to reach Category 4 and 5 strength. A 2019 study indicates that climate change has been driving the observed trend of rapid intensification of tropical cyclones in the Atlantic basin. Rapidly intensifying cyclones are hard to forecast and therefore pose additional risk to coastal communities.
Warmer air can hold more water vapor: the theoretical maximum water vapor content is given by the Clausius–Clapeyron relation, which yields ≈7% increase in water vapor in the atmosphere per warming. All models that were assessed in a 2019 review paper show a future increase of rainfall rates. Additional sea level rise will increase storm surge levels. It is plausible that extreme wind waves see an increase as a consequence of changes in tropical cyclones, further exacerbating storm surge dangers to coastal communities. The compounding effects from floods, storm surge, and terrestrial flooding (rivers) are projected to increase due to global warming.
There is currently no consensus on how climate change will affect the overall frequency of tropical cyclones. A majority of climate models show a decreased frequency in future projections. For instance, a 2020 paper comparing nine high-resolution climate models found robust decreases in frequency in the Southern Indian Ocean and the Southern Hemisphere more generally, while finding mixed signals for Northern Hemisphere tropical cyclones. Observations have shown little change in the overall frequency of tropical cyclones worldwide, with increased frequency in the North Atlantic and central Pacific, and significant decreases in the southern Indian Ocean and western North Pacific.
There has been a poleward expansion of the latitude at which the maximum intensity of tropical cyclones occurs, which may be associated with climate change. In the North Pacific, there may also have been an eastward expansion. Between 1949 and 2016, there was a slowdown in tropical cyclone translation speeds. It is unclear still to what extent this can be attributed to climate change: climate models do not all show this feature.
A 2021 study review article concluded that the geographic range of tropical cyclones will probably expand poleward in response to climate warming of the Hadley circulation.
When hurricane winds speed rise by 5%, its destructive power rise by about 50%. Therefore, as climate change increased the wind speed of Hurricane Helene by 11%, it increased the destruction from it by more than twice. According to World Weather Attribution the influence of climate change on the rainfall of some latest hurricanes can be described as follows:
Intensity
Tropical cyclone intensity is based on wind speeds and pressure. Relationships between winds and pressure are often used in determining the intensity of a storm. Tropical cyclone scales, such as the Saffir-Simpson hurricane wind scale and Australia's scale (Bureau of Meteorology), only use wind speed for determining the category of a storm. The most intense storm on record is Typhoon Tip in the northwestern Pacific Ocean in 1979, which reached a minimum pressure of and maximum sustained wind speeds of . The highest maximum sustained wind speed ever recorded was in Hurricane Patricia in 2015—the most intense cyclone ever recorded in the Western Hemisphere.
Factors that influence intensity
Warm sea surface temperatures are required for tropical cyclones to form and strengthen. The commonly-accepted minimum temperature range for this to occur is , however, multiple studies have proposed a lower minimum of . Higher sea surface temperatures result in faster intensification rates and sometimes even rapid intensification. High ocean heat content, also known as Tropical Cyclone Heat Potential, allows storms to achieve a higher intensity. Most tropical cyclones that experience rapid intensification are traversing regions of high ocean heat content rather than lower values. High ocean heat content values can help to offset the oceanic cooling caused by the passage of a tropical cyclone, limiting the effect this cooling has on the storm. Faster-moving systems are able to intensify to higher intensities with lower ocean heat content values. Slower-moving systems require higher values of ocean heat content to achieve the same intensity.
The passage of a tropical cyclone over the ocean causes the upper layers of the ocean to cool substantially, a process known as upwelling, which can negatively influence subsequent cyclone development. This cooling is primarily caused by wind-driven mixing of cold water from deeper in the ocean with the warm surface waters. This effect results in a negative feedback process that can inhibit further development or lead to weakening. Additional cooling may come in the form of cold water from falling raindrops (this is because the atmosphere is cooler at higher altitudes). Cloud cover may also play a role in cooling the ocean, by shielding the ocean surface from direct sunlight before and slightly after the storm passage. All these effects can combine to produce a dramatic drop in sea surface temperature over a large area in just a few days. Conversely, the mixing of the sea can result in heat being inserted in deeper waters, with potential effects on global climate.
Vertical wind shear decreases tropical cyclone predicability, with storms exhibiting wide range of responses in the presence of shear. Wind shear often negatively affects tropical cyclone intensification by displacing moisture and heat from a system's center. Low levels of vertical wind shear are most optimal for strengthening, while stronger wind shear induces weakening. Dry air entraining into a tropical cyclone's core has a negative effect on its development and intensity by diminishing atmospheric convection and introducing asymmetries in the storm's structure. Symmetric, strong outflow leads to a faster rate of intensification than observed in other systems by mitigating local wind shear. Weakening outflow is associated with the weakening of rainbands within a tropical cyclone. Tropical cyclones may still intensify, even rapidly, in the presence of moderate or strong wind shear depending on the evolution and structure of the storm's convection.
The size of tropical cyclones plays a role in how quickly they intensify. Smaller tropical cyclones are more prone to rapid intensification than larger ones. The Fujiwhara effect, which involves interaction between two tropical cyclones, can weaken and ultimately result in the dissipation of the weaker of two tropical cyclones by reducing the organization of the system's convection and imparting horizontal wind shear. Tropical cyclones typically weaken while situated over a landmass because conditions are often unfavorable as a result of the lack of oceanic forcing. The Brown ocean effect can allow a tropical cyclone to maintain or increase its intensity following landfall, in cases where there has been copious rainfall, through the release of latent heat from the saturated soil. Orographic lift can cause a significant increase in the intensity of the convection of a tropical cyclone when its eye moves over a mountain, breaking the capped boundary layer that had been restraining it. Jet streams can both enhance and inhibit tropical cyclone intensity by influencing the storm's outflow as well as vertical wind shear.
Rapid intensification
On occasion, tropical cyclones may undergo a process known as rapid intensification, a period in which the maximum sustained winds of a tropical cyclone increase by or more within 24 hours. Similarly, rapid deepening in tropical cyclones is defined as a minimum sea surface pressure decrease of per hour or within a 24-hour period; explosive deepening occurs when the surface pressure decreases by per hour for at least 12 hours or per hour for at least 6 hours.
For rapid intensification to occur, several conditions must be in place. Water temperatures must be extremely high, near or above , and water of this temperature must be sufficiently deep such that waves do not upwell cooler waters to the surface. On the other hand, Tropical Cyclone Heat Potential is one of such non-conventional subsurface oceanographic parameters influencing the cyclone intensity.
Wind shear must be low. When wind shear is high, the convection and circulation in the cyclone will be disrupted. Usually, an anticyclone in the upper layers of the troposphere above the storm must be present as well—for extremely low surface pressures to develop, air must be rising very rapidly in the eyewall of the storm, and an upper-level anticyclone helps channel this air away from the cyclone efficiently. However, some cyclones such as Hurricane Epsilon have rapidly intensified despite relatively unfavorable conditions.
Dissipation
There are a number of ways a tropical cyclone can weaken, dissipate, or lose its tropical characteristics. These include making landfall, moving over cooler water, encountering dry air, or interacting with other weather systems; however, once a system has dissipated or lost its tropical characteristics, its remnants could regenerate a tropical cyclone if environmental conditions become favorable.
A tropical cyclone can dissipate when it moves over waters significantly cooler than . This will deprive the storm of such tropical characteristics as a warm core with thunderstorms near the center, so that it becomes a remnant low-pressure area. Remnant systems may persist for several days before losing their identity. This dissipation mechanism is most common in the eastern North Pacific. Weakening or dissipation can also occur if a storm experiences vertical wind shear which causes the convection and heat engine to move away from the center. This normally ceases the development of a tropical cyclone. In addition, its interaction with the main belt of the Westerlies, by means of merging with a nearby frontal zone, can cause tropical cyclones to evolve into extratropical cyclones. This transition can take 1–3 days.
Should a tropical cyclone make landfall or pass over an island, its circulation could start to break down, especially if it encounters mountainous terrain. When a system makes landfall on a large landmass, it is cut off from its supply of warm moist maritime air and starts to draw in dry continental air. This, combined with the increased friction over land areas, leads to the weakening and dissipation of the tropical cyclone. Over a mountainous terrain, a system can quickly weaken. Over flat areas, it may endure for two to three days before circulation breaks down and dissipates.
Over the years, there have been a number of techniques considered to try to artificially modify tropical cyclones. These techniques have included using nuclear weapons, cooling the ocean with icebergs, blowing the storm away from land with giant fans, and seeding selected storms with dry ice or silver iodide. These techniques, however, fail to appreciate the duration, intensity, power or size of tropical cyclones.
Methods for assessing intensity
A variety of methods or techniques, including surface, satellite, and aerial, are used to assess the intensity of a tropical cyclone. Reconnaissance aircraft fly around and through tropical cyclones, outfitted with specialized instruments, to collect information that can be used to ascertain the winds and pressure of a system. Tropical cyclones possess winds of different speeds at different heights. Winds recorded at flight level can be converted to find the wind speeds at the surface. Surface observations, such as ship reports, land stations, mesonets, coastal stations, and buoys, can provide information on a tropical cyclone's intensity or the direction it is traveling.
Wind-pressure relationships (WPRs) are used as a way to determine the pressure of a storm based on its wind speed. Several different methods and equations have been proposed to calculate WPRs. Tropical cyclones agencies each use their own, fixed WPR, which can result in inaccuracies between agencies that are issuing estimates on the same system. The ASCAT is a scatterometer used by the MetOp satellites to map the wind field vectors of tropical cyclones. The SMAP uses an L-band radiometer channel to determine the wind speeds of tropical cyclones at the ocean surface, and has been shown to be reliable at higher intensities and under heavy rainfall conditions, unlike scatterometer-based and other radiometer-based instruments.
The Dvorak technique plays a large role in both the classification of a tropical cyclone and the determination of its intensity. Used in warning centers, the method was developed by Vernon Dvorak in the 1970s, and uses both visible and infrared satellite imagery in the assessment of tropical cyclone intensity. The Dvorak technique uses a scale of "T-numbers", scaling in increments of 0.5 from T1.0 to T8.0. Each T-number has an intensity assigned to it, with larger T-numbers indicating a stronger system. Tropical cyclones are assessed by forecasters according to an array of patterns, including curved banding features, shear, central dense overcast, and eye, to determine the T-number and thus assess the intensity of the storm.
The Cooperative Institute for Meteorological Satellite Studies works to develop and improve automated satellite methods, such as the Advanced Dvorak Technique (ADT) and SATCON. The ADT, used by a large number of forecasting centers, uses infrared geostationary satellite imagery and an algorithm based upon the Dvorak technique to assess the intensity of tropical cyclones. The ADT has a number of differences from the conventional Dvorak technique, including changes to intensity constraint rules and the usage of microwave imagery to base a system's intensity upon its internal structure, which prevents the intensity from leveling off before an eye emerges in infrared imagery. The SATCON weights estimates from various satellite-based systems and microwave sounders, accounting for the strengths and flaws in each individual estimate, to produce a consensus estimate of a tropical cyclone's intensity which can be more reliable than the Dvorak technique at times.
Intensity metrics
Multiple intensity metrics are used, including accumulated cyclone energy (ACE), the Hurricane Surge Index, the Hurricane Severity Index, the Power Dissipation Index (PDI), and integrated kinetic energy (IKE). ACE is a metric of the total energy a system has exerted over its lifespan. ACE is calculated by summing the squares of a cyclone's sustained wind speed, every six hours as long as the system is at or above tropical storm intensity and either tropical or subtropical. The calculation of the PDI is similar in nature to ACE, with the major difference being that wind speeds are cubed rather than squared.
The Hurricane Surge Index is a metric of the potential damage a storm may inflict via storm surge. It is calculated by squaring the dividend of the storm's wind speed and a climatological value (), and then multiplying that quantity by the dividend of the radius of hurricane-force winds and its climatological value (). This can be represented in equation form as:
where is the storm's wind speed and is the radius of hurricane-force winds. The Hurricane Severity Index is a scale that can assign up to 50 points to a system; up to 25 points come from intensity, while the other 25 come from the size of the storm's wind field. The IKE model measures the destructive capability of a tropical cyclone via winds, waves, and surge. It is calculated as:
where is the density of air, is a sustained surface wind speed value, and is the volume element.
Classification and naming
Classification
Around the world, tropical cyclones are classified in different ways, based on the location (tropical cyclone basins), the structure of the system and its intensity. For example, within the Northern Atlantic and Eastern Pacific basins, a tropical cyclone with wind speeds of over is called a hurricane, while it is called a typhoon or a severe cyclonic storm within the Western Pacific or North Indian oceans. When a hurricane passes west across the International Dateline in the Northern Hemisphere, it becomes known as a typhoon. This happened in 2014 for Hurricane Genevieve, which became Typhoon Genevieve.
Within the Southern Hemisphere, it is either called a hurricane, tropical cyclone or a severe tropical cyclone, depending on if it is located within the South Atlantic, South-West Indian Ocean, Australian region or the South Pacific Ocean. The descriptors for tropical cyclones with wind speeds below vary by tropical cyclone basin and may be further subdivided into categories such as "tropical storm", "cyclonic storm", "tropical depression", or "deep depression".
Naming
The practice of using given names to identify tropical cyclones dates back to the late 1800s and early 1900s and gradually superseded the existing system—simply naming cyclones based on what they hit. The system currently used provides positive identification of severe weather systems in a brief form, that is readily understood and recognized by the public. The credit for the first usage of personal names for weather systems is generally given to the Queensland Government Meteorologist Clement Wragge who named systems between 1887 and 1907. This system of naming weather systems fell into disuse for several years after Wragge retired, until it was revived in the latter part of World War II for the Western Pacific. Formal naming schemes have subsequently been introduced for the North and South Atlantic, Eastern, Central, Western and Southern Pacific basins as well as the Australian region and Indian Ocean.
At present, tropical cyclones are officially named by one of twelve meteorological services and retain their names throughout their lifetimes to provide ease of communication between forecasters and the general public regarding forecasts, watches, and warnings. Since the systems can last a week or longer, and more than one can be occurring in the same basin at the same time, the names are thought to reduce the confusion about what storm is being described. Names are assigned in order from predetermined lists with one, three, or ten-minute sustained wind speeds of more than depending on which basin it originates.
Standards vary from basin to basin. Some tropical depressions are named in the Western Pacific. Tropical cyclones have to have a significant amount of gale-force winds occurring around the center before they are named within the Southern Hemisphere. The names of significant tropical cyclones in the North Atlantic Ocean, Pacific Ocean, and Australian region are retired from the naming lists and replaced with another name. Tropical cyclones that develop around the world are assigned an identification code consisting of a two-digit number and suffix letter by the warning centers that monitor them.
Related cyclone types
In addition to tropical cyclones, there are two other classes of cyclones within the spectrum of cyclone types. These kinds of cyclones, known as extratropical cyclones and subtropical cyclones, can be stages a tropical cyclone passes through during its formation or dissipation. An extratropical cyclone is a storm that derives energy from horizontal temperature differences, which are typical in higher latitudes. A tropical cyclone can become extratropical as it moves toward higher latitudes if its energy source changes from heat released by condensation to differences in temperature between air masses. Although not as frequently, an extratropical cyclone can transform into a subtropical storm, and from there into a tropical cyclone. From space, extratropical storms have a characteristic "comma-shaped" cloud pattern. Extratropical cyclones can also be dangerous when their low-pressure centers cause powerful winds and high seas.
A subtropical cyclone is a weather system that has some characteristics of a tropical cyclone and some characteristics of an extratropical cyclone. They can form in a wide band of latitudes, from the equator to 50°. Although subtropical storms rarely have hurricane-force winds, they may become tropical in nature as their cores warm.
Structure
Eye and center
At the center of a mature tropical cyclone, air sinks rather than rises. For a sufficiently strong storm, air may sink over a layer deep enough to suppress cloud formation, thereby creating a clear "eye". Weather in the eye is normally calm and free of convective clouds, although the sea may be extremely violent. The eye is normally circular and is typically in diameter, though eyes as small as and as large as have been observed.
The cloudy outer edge of the eye is called the "eyewall". The eyewall typically expands outward with height, resembling an arena football stadium; this phenomenon is sometimes referred to as the "stadium effect". The eyewall is where the greatest wind speeds are found, air rises most rapidly, clouds reach their highest altitude, and precipitation is the heaviest. The heaviest wind damage occurs where a tropical cyclone's eyewall passes over land.
In a weaker storm, the eye may be obscured by the central dense overcast, which is the upper-level cirrus shield that is associated with a concentrated area of strong thunderstorm activity near the center of a tropical cyclone.
The eyewall may vary over time in the form of eyewall replacement cycles, particularly in intense tropical cyclones. Outer rainbands can organize into an outer ring of thunderstorms that slowly moves inward, which is believed to rob the primary eyewall of moisture and angular momentum. When the primary eyewall weakens, the tropical cyclone weakens temporarily. The outer eyewall eventually replaces the primary one at the end of the cycle, at which time the storm may return to its original intensity.
Size
There are a variety of metrics commonly used to measure storm size. The most common metrics include the radius of maximum wind, the radius of wind (i.e. gale force), the radius of outermost closed isobar (ROCI), and the radius of vanishing wind. An additional metric is the radius at which the cyclone's relative vorticity field decreases to 1×10−5 s−1.
On Earth, tropical cyclones span a large range of sizes, from as measured by the radius of vanishing wind. They are largest on average in the northwest Pacific Ocean basin and smallest in the northeastern Pacific Ocean basin. If the radius of outermost closed isobar is less than two degrees of latitude (), then the cyclone is "very small" or a "midget". A radius of 3–6 latitude degrees () is considered "average sized". "Very large" tropical cyclones have a radius of greater than 8 degrees (). Observations indicate that size is only weakly correlated to variables such as storm intensity (i.e. maximum wind speed), radius of maximum wind, latitude, and maximum potential intensity. Typhoon Tip is the largest cyclone on record, with tropical storm-force winds in diameter. The smallest storm on record is Tropical Storm Marco of 2008, which generated tropical storm-force winds only in diameter.
Movement
The movement of a tropical cyclone (i.e. its "track") is typically approximated as the sum of two terms: "steering" by the background environmental wind and "beta drift". Some tropical cyclones can move across large distances, such as Hurricane John, the second longest-lasting tropical cyclone on record, which traveled , the longest track of any Northern Hemisphere tropical cyclone, over its 31-day lifespan in 1994.
Environmental steering
Environmental steering is the primary influence on the motion of tropical cyclones. It represents the movement of the storm due to prevailing winds and other wider environmental conditions, similar to "leaves carried along by a stream".
Physically, the winds, or flow field, in the vicinity of a tropical cyclone may be treated as having two parts: the flow associated with the storm itself, and the large-scale background flow of the environment. Tropical cyclones can be treated as local maxima of vorticity suspended within the large-scale background flow of the environment. In this way, tropical cyclone motion may be represented to first-order as advection of the storm by the local environmental flow. This environmental flow is termed the "steering flow" and is the dominant influence on tropical cyclone motion. The strength and direction of the steering flow can be approximated as a vertical integration of the winds blowing horizontally in the cyclone's vicinity, weighted by the altitude at which those winds are occurring. Because winds can vary with height, determining the steering flow precisely can be difficult.
The pressure altitude at which the background winds are most correlated with a tropical cyclone's motion is known as the "steering level". The motion of stronger tropical cyclones is more correlated with the background flow averaged across a thicker portion of troposphere compared to weaker tropical cyclones whose motion is more correlated with the background flow averaged across a narrower extent of the lower troposphere. When wind shear and latent heat release is present, tropical cyclones tend to move towards regions where potential vorticity is increasing most quickly.
Climatologically, tropical cyclones are steered primarily westward by the east-to-west trade winds on the equatorial side of the subtropical ridge—a persistent high-pressure area over the world's subtropical oceans. In the tropical North Atlantic and Northeast Pacific oceans, the trade winds steer tropical easterly waves westward from the African coast toward the Caribbean Sea, North America, and ultimately into the central Pacific Ocean before the waves dampen out. These waves are the precursors to many tropical cyclones within this region. In contrast, in the Indian Ocean and Western Pacific in both hemispheres, tropical cyclogenesis is influenced less by tropical easterly waves and more by the seasonal movement of the Intertropical Convergence Zone and the monsoon trough. Other weather systems such as mid-latitude troughs and broad monsoon gyres can also influence tropical cyclone motion by modifying the steering flow.
Beta drift
In addition to environmental steering, a tropical cyclone will tend to drift poleward and westward, a motion known as "beta drift". This motion is due to the superposition of a vortex, such as a tropical cyclone, onto an environment in which the Coriolis force varies with latitude, such as on a sphere or beta plane. The magnitude of the component of tropical cyclone motion associated with the beta drift ranges between and tends to be larger for more intense tropical cyclones and at higher latitudes. It is induced indirectly by the storm itself as a result of feedback between the cyclonic flow of the storm and its environment.
Physically, the cyclonic circulation of the storm advects environmental air poleward east of center and equatorial west of center. Because air must conserve its angular momentum, this flow configuration induces a cyclonic gyre equatorward and westward of the storm center and an anticyclonic gyre poleward and eastward of the storm center. The combined flow of these gyres acts to advect the storm slowly poleward and westward. This effect occurs even if there is zero environmental flow. Due to a direct dependence of the beta drift on angular momentum, the size of a tropical cyclone can affect the influence of beta drift on its motion; beta drift imparts a greater influence on the movement of larger tropical cyclones than that of smaller ones.
Multiple storm interaction
A third component of motion that occurs relatively infrequently involves the interaction of multiple tropical cyclones. When two cyclones approach one another, their centers will begin orbiting cyclonically about a point between the two systems. Depending on their separation distance and strength, the two vortices may simply orbit around one another, or else may spiral into the center point and merge. When the two vortices are of unequal size, the larger vortex will tend to dominate the interaction, and the smaller vortex will orbit around it. This phenomenon is called the Fujiwhara effect, after Sakuhei Fujiwhara.
Interaction with the mid-latitude westerlies
Though a tropical cyclone typically moves from east to west in the tropics, its track may shift poleward and eastward either as it moves west of the subtropical ridge axis or else if it interacts with the mid-latitude flow, such as the jet stream or an extratropical cyclone. This motion, termed "recurvature", commonly occurs near the western edge of the major ocean basins, where the jet stream typically has a poleward component and extratropical cyclones are common. An example of tropical cyclone recurvature was Typhoon Ioke in 2006.
Effects
Natural phenomena caused or worsened by tropical cyclones
Tropical cyclones out at sea cause large waves, heavy rain, floods and high winds, disrupting international shipping and, at times, causing shipwrecks. Tropical cyclones stir up water, leaving a cool wake behind them, which causes the region to be less favorable for subsequent tropical cyclones. On land, strong winds can damage or destroy vehicles, buildings, bridges, and other outside objects, turning loose debris into deadly flying projectiles. The storm surge, or the increase in sea level due to the cyclone, is typically the worst effect from landfalling tropical cyclones, historically resulting in 90% of tropical cyclone deaths. Cyclone Mahina produced the highest storm surge on record, , at Bathurst Bay, Queensland, Australia, in March 1899.
Other ocean-based hazards that tropical cyclones produce are rip currents and undertow. These hazards can occur hundreds of kilometers (hundreds of miles) away from the center of a cyclone, even if other weather conditions are favorable.
The broad rotation of a landfalling tropical cyclone, and vertical wind shear at its periphery, spawns tornadoes. Tornadoes can also be spawned as a result of eyewall mesovortices, which persist until landfall. Hurricane Ivan produced 120 tornadoes, more than any other tropical cyclone. Lightning activity is produced within tropical cyclones. This activity is more intense within stronger storms and closer to and within the storm's eyewall. Tropical cyclones can increase the amount of snowfall a region experiences by delivering additional moisture. Wildfires can be worsened when a nearby storm fans their flames with its strong winds.
Effect on property and human life
Tropical cyclones regularly affect the coastlines of most of Earth's major bodies of water along the Atlantic, Pacific, and Indian oceans. Tropical cyclones have caused significant destruction and loss of human life, resulting in about 2 million deaths since the 19th century. Large areas of standing water caused by flooding lead to infection, as well as contributing to mosquito-borne illnesses. Crowded evacuees in shelters increase the risk of disease propagation. Tropical cyclones significantly interrupt infrastructure, leading to power outages, bridge and road destruction, and the hampering of reconstruction efforts.
Winds and water from storms can damage or destroy homes, buildings, and other manmade structures. Tropical cyclones destroy agriculture, kill livestock, and prevent access to marketplaces for both buyers and sellers; both of these result in financial losses. Powerful cyclones that make landfall – moving from the ocean to over land – are some of the most powerful, although that is not always the case. An average of 86 tropical cyclones of tropical storm intensity form annually worldwide, with 47 reaching hurricane or typhoon strength, and 20 becoming intense tropical cyclones, super typhoons, or major hurricanes (at least of Category 3 intensity).
Africa
In Africa, tropical cyclones can originate from tropical waves generated over the Sahara Desert, or otherwise strike the Horn of Africa and Southern Africa. Cyclone Idai in March 2019 hit central Mozambique, becoming the deadliest tropical cyclone on record in Africa, with 1,302 fatalities, and damage estimated at US$2.2 billion. Réunion island, located east of Southern Africa, experiences some of the wettest tropical cyclones on record. In January 1980, Cyclone Hyacinthe produced 6,083 mm (239.5 in) of rain over 15 days, which was the largest rain total recorded from a tropical cyclone on record.
Asia
In Asia, tropical cyclones from the Indian and Pacific oceans regularly affect some of the most populated countries on Earth. In 1970, a cyclone struck Bangladesh, then known as East Pakistan, producing a storm surge that killed at least 300,000 people. This made it the deadliest tropical cyclone on record. In October 2019, Typhoon Hagibis struck the Japanese island of Honshu and inflicted US$15 billion in damage, making it the costliest storm on record in Japan. The islands that comprise Oceania, from Australia to French Polynesia, are routinely affected by tropical cyclones. In Indonesia, a cyclone struck the island of Flores in April 1973, killing 1,653 people, making it the deadliest tropical cyclone recorded in the Southern Hemisphere.
North and South America
Atlantic and Pacific hurricanes regularly affect North America. In the United States, hurricanes Katrina in 2005 and Harvey in 2017 are the country's costliest ever natural disasters, with monetary damage estimated at US$125 billion. Katrina struck Louisiana and largely destroyed the city of New Orleans, while Harvey caused significant flooding in southeastern Texas after it dropped of rainfall; this was the highest rainfall total on record in the country.
The Caribbean islands are regularly hit by hurricanes, which have caused multiple humanitarian crises in Haiti since 2004 due in part to the lack of infrastructure and high population density in urban areas. In 2004, hurricane Jeanne caused severe flooding and mudslides, and a total estimated 3,006 deaths. More recently, in 2016, hurricane Matthew caused US$2.8 billion in damages, killing an estimated 674 people.
The northern portion of South America experiences occasional tropical cyclones, with 173 fatalities from Tropical Storm Bret in August 1993. The South Atlantic Ocean is generally inhospitable to the formation of a tropical storm. However, in March 2004, Hurricane Catarina struck southeastern Brazil as the first hurricane on record in the South Atlantic Ocean.
Europe
Europe is rarely affected by tropical cyclones; however, the continent regularly encounters storms after they transitioned into extratropical cyclones. Only one tropical depression – Vince in 2005 – struck Spain, and only one subtropical cyclone – Subtropical Storm Alpha in 2020 – struck Portugal. Occasionally, there are tropical-like cyclones in the Mediterranean Sea.
Environmental effects
Although cyclones take an enormous toll in lives and personal property, they may be important factors in the precipitation regimes of places they affect, as they may bring much-needed precipitation to otherwise dry regions. Their precipitation may also alleviate drought conditions by restoring soil moisture, though one study focused on the Southeastern United States suggested tropical cyclones did not offer significant drought recovery. Tropical cyclones also help maintain the global heat balance by moving warm, moist tropical air to the middle latitudes and polar regions, and by regulating the thermohaline circulation through upwelling. Research on Pacific cyclones has demonstrated that deeper layers of the ocean receive a heat transfer from these powerful storms.
The storm surge and winds of hurricanes may be destructive to human-made structures, but they also stir up the waters of coastal estuaries, which are typically important fish breeding locales. Ecosystems, such as saltmarshes and Mangrove forests, can be severely damaged or destroyed by tropical cyclones, which erode land and destroy vegetation. Tropical cyclones can cause harmful algae blooms to form in bodies of water by increasing the amount of nutrients available. Insect populations can decrease in both quantity and diversity after the passage of storms. Strong winds associated with tropical cyclones and their remnants are capable of felling thousands of trees, causing damage to forests.
When hurricanes surge upon shore from the ocean, salt is introduced to many freshwater areas and raises the salinity levels too high for some habitats to withstand. Some are able to cope with the salt and recycle it back into the ocean, but others can not release the extra surface water quickly enough or do not have a large enough freshwater source to replace it. Because of this, some species of plants and vegetation die due to the excess salt. Hurricanes can carry toxins and acids onshore when they make landfall. The floodwater can pick up the toxins from different spills and contaminate the land that it passes over. These toxins are harmful to the people and animals in the area, as well as the environment around them. Tropical cyclones can cause oil spills by damaging or destroying pipelines and storage facilities. Similarly, chemical spills have been reported when chemical and processing facilities were damaged. Waterways have become contaminated with toxic levels of metals such as nickel, chromium, and mercury during tropical cyclones.
Tropical cyclones can have an extensive effect on geography, such as creating or destroying land. Cyclone Bebe increased the size of Tuvalu island, Funafuti Atoll, by nearly 20%. Hurricane Walaka destroyed the small East Island in 2018, which destroyed the habitat for the endangered Hawaiian monk seal, as well as, threatened sea turtles and seabirds. Landslides frequently occur during tropical cyclones and can vastly alter landscapes. Some storms are capable of causing hundreds to tens of thousands of landslides. Storms can erode coastlines over an extensive area and transport the sediment to other locations.
Observation and forecasting
Observation
Tropical cyclones have occurred around the world for millennia. Reanalyses and research are being undertaken to extend the historical record, through the usage of proxy data such as overwash deposits, beach ridges and historical documents such as diaries. Major tropical cyclones leave traces in overwash records and shell layers in some coastal areas, which have been used to gain insight into hurricane activity over the past thousands of years. Sediment records in Western Australia suggest an intense tropical cyclone in the 4th millennium BC.
Proxy records based on paleotempestological research have revealed that major hurricane activity along the Gulf of Mexico coast varies on timescales of centuries to millennia. In the year 957, a powerful typhoon struck southern China, killing around 10,000 people due to flooding. The Spanish colonization of Mexico described "tempestades" in 1730, although the official record for Pacific hurricanes only dates to 1949. In the south-west Indian Ocean, the tropical cyclone record goes back to 1848. In 2003, the Atlantic hurricane reanalysis project examined and analyzed the historical record of tropical cyclones in the Atlantic back to 1851, extending the existing database from 1886.
Before satellite imagery became available during the 20th century, many of these systems went undetected unless it impacted land or a ship encountered it by chance. Often in part because of the threat of hurricanes, many coastal regions had sparse population between major ports until the advent of automobile tourism; therefore, the most severe portions of hurricanes striking the coast may have gone unmeasured in some instances. The combined effects of ship destruction and remote landfall severely limit the number of intense hurricanes in the official record before the era of hurricane reconnaissance aircraft and satellite meteorology. Although the record shows a distinct increase in the number and strength of intense hurricanes, therefore, experts regard the early data as suspect. The ability of climatologists to make a long-term analysis of tropical cyclones is limited by the amount of reliable historical data.
In the 1940s, routine aircraft reconnaissance started in both the Atlantic and Western Pacific basin in the mid-1940s, which provided ground truth data. Early flights were only made once or twice a day. In 1960, Polar-orbiting weather satellites were first launched by the United States National Aeronautics and Space Administration, but were not declared operational until 1965. It took several years for some of the warning centers to take advantage of this new viewing platform and develop the expertise to associate satellite signatures with storm position and intensity.
Intense tropical cyclones pose a particular observation challenge, as they are a dangerous oceanic phenomenon, and weather stations, being relatively sparse, are rarely available on the site of the storm itself. In general, surface observations are available only if the storm is passing over an island or a coastal area, or if there is a nearby ship. Real-time measurements are usually taken in the periphery of the cyclone, where conditions are less catastrophic and its true strength cannot be evaluated. For this reason, there are teams of meteorologists that move into the path of tropical cyclones to help evaluate their strength at the point of landfall.
Tropical cyclones are tracked by weather satellites capturing visible and infrared images from space, usually at half-hour to quarter-hour intervals. As a storm approaches land, it can be observed by land-based Doppler weather radar. Radar plays a crucial role around landfall by showing a storm's location and intensity every several minutes. Other satellites provide information from the perturbations of GPS signals, providing thousands of snapshots per day and capturing atmospheric temperature, pressure, and moisture content.
In situ measurements, in real-time, can be taken by sending specially equipped reconnaissance flights into the cyclone. In the Atlantic basin, these flights are regularly flown by United States government hurricane hunters. These aircraft fly directly into the cyclone and take direct and remote-sensing measurements. The aircraft launch GPS dropsondes inside the cyclone. These sondes measure temperature, humidity, pressure, and especially winds between flight level and the ocean's surface. A new era in hurricane observation began when a remotely piloted Aerosonde, a small drone aircraft, was flown through Tropical Storm Ophelia as it passed Virginia's eastern shore during the 2005 hurricane season. A similar mission was also completed successfully in the western Pacific Ocean.
Forecasting
High-speed computers and sophisticated simulation software allow forecasters to produce computer models that predict tropical cyclone tracks based on the future position and strength of high- and low-pressure systems. Combining forecast models with increased understanding of the forces that act on tropical cyclones, as well as with a wealth of data from Earth-orbiting satellites and other sensors, scientists have increased the accuracy of track forecasts over recent decades.
However, scientists are not as skillful at predicting the intensity of tropical cyclones. The lack of improvement in intensity forecasting is attributed to the complexity of tropical systems and an incomplete understanding of factors that affect their development. New tropical cyclone position and forecast information is available at least every six hours from the various warning centers.
Geopotential height
In meteorology, geopotential heights are used when creating forecasts and analyzing pressure systems. Geopotential heights represent the estimate of the real height of a pressure system above the average sea level. Geopotential heights for weather are divided up into several levels. The lowest geopotential height level is , which represents the lowest of the atmosphere. The moisture content, gained by using either the relative humidity or the precipitable water value, is used in creating forecasts for precipitation.
The next level, , is at a height of . 700 hPa is regarded as the highest point in the lower atmosphere. At this layer, both vertical movement and moisture levels are used to locate and create forecasts for precipitation. The middle level of the atmosphere is at or a height of . The 500 hPa level is used for measuring atmospheric vorticity, commonly known as the spin of air. The relative humidity is also analyzed at this height to establish where precipitation is likely to materialize. The next level occurs at or a height of . The top-most level is located at , which corresponds to a height of . Both the 200 and 300 hPa levels are mainly used to locate the jet stream.
Society and culture
Preparations
Ahead of the formal season starting, people are urged to prepare for the effects of a tropical cyclone by politicians and weather forecasters, among others. They prepare by determining their risk to the different types of weather, tropical cyclones cause, checking their insurance coverage and emergency supplies, as well as determining where to evacuate to if needed. When a tropical cyclone develops and is forecast to impact land, each member nation of the World Meteorological Organization issues various watches and warnings to cover the expected effects. However, there are some exceptions with the United States National Hurricane Center and Fiji Meteorological Service responsible for issuing or recommending warnings for other nations in their area of responsibility.
An important decision in individual preparedness is determining if and when to evacuate an area that will be affected by a tropical cyclone. Tropical cyclone tracking charts allow people to track ongoing systems to form their own opinions regarding where the storms are going and whether or not they need to prepare for the system being tracked, including possible evacuation. This continues to be encouraged by the National Oceanic and Atmospheric Administration and National Hurricane Center.
Response
Hurricane response is the disaster response after a hurricane. Activities performed by hurricane responders include assessment, restoration, and demolition of buildings; removal of debris and waste; repairs to land-based and maritime infrastructure; and public health services including search and rescue operations. Hurricane response requires coordination between federal, tribal, state, local, and private entities. According to the National Voluntary Organizations Active in Disaster, potential response volunteers should affiliate with established organizations and should not self-deploy, so that proper training and support can be provided to mitigate the danger and stress of response work.
Hurricane responders face many hazards. Hurricane responders may be exposed to chemical and biological contaminants including stored chemicals, sewage, human remains, and mold growth encouraged by flooding, as well as asbestos and lead that may be present in older buildings. Common injuries arise from falls from heights, such as from a ladder or from level surfaces; from electrocution in flooded areas, including from backfeed from portable generators; or from motor vehicle accidents. Long and irregular shifts may lead to sleep deprivation and fatigue, increasing the risk of injuries, and workers may experience mental stress associated with a traumatic incident. Heat stress is a concern as workers are often exposed to hot and humid temperatures, wear protective clothing and equipment, and have physically difficult tasks.
| Physical sciences | Earth science | null |
26734587 | https://en.wikipedia.org/wiki/Orthotics | Orthotics | Orthotics () is a medical specialty that focuses on the design and application of orthoses, sometimes known as braces, calipers, or splints. An is "an externally applied device used to influence the structural and functional characteristics of the neuromuscular and skeletal systems." Orthotists are medical professionals who specialize in designing orthotic devices such as braces or foot orthoses.
Classification
Orthotic devices are classified into four areas of the body according to the international classification system (ICS): orthotics of the lower extremities, orthotics of the upper extremities, orthotics for the trunk, and orthotics for the head. Orthoses are also classified by function: paralysis orthoses and relief orthoses.
Under the International Standard terminology, orthoses are classified by an acronym describing the anatomical joints they support. Some examples include KAFO, or knee-ankle-foot orthoses, which span the knee, ankle, and foot; TLSO, or thoracic-lumbar-sacral orthoses, supporting the thoracic, lumbar and sacral regions of the spine. The use of the International Standard is promoted to reduce the widespread variation in the description of orthoses, which is often a barrier to interpreting research studies.
The transition from an orthosis to a prosthesis can be fluid. An example is compensating for a leg length discrepancy, equivalent to replacing a missing part of a limb. Another example is the replacement of the forefoot after a forefoot amputation. This treatment is often made from a combination of a prosthesis to replace the forefoot and an orthosis to replace the lost muscular function (ortho prosthesis).
Orthotist
An orthotist is a specialist responsible for the customising, manufacture, and repair of orthotic devices (orthoses). The manufacture of modern orthoses requires both artistic skills in modeling body shapes and manual skills in processing traditional and innovative materials— CAD/CAM, CNC machines and 3D printing are involved in orthotic manufacture. Orthotics also combines knowledge of anatomy and physiology, pathophysiology, biomechanics and engineering.
In the United States, while orthotists require a prescription from a licensed healthcare provider, physical therapists are not legally authorized to prescribe orthoses. In the U.K., orthotists will often accept referrals from doctors or other healthcare professionals for orthotic assessment without requiring a prescription.
Prescription and manufacturing
Orthoses are offered as:
Custom-fabricated products – they are in the foreground of an optimal supply and are individually manufactured. If the physical examination of a patient is carried out precisely, the clinical picture often shows a combination of several functional deviations. Each functional deviation can be slight or severe. The combination of the functional deviation and its characteristics leads to a detailed indication. A major advantage of custom-made products is that the various necessary orthotic functions when doing the configuration of the orthotics can be optimally matched to the determined functional deviations. Another advantage of custom-made products is that each orthosis is made to fit the individual body shape of the patient. Custom-fabricated products were traditionally made by following a trace of the extremity with measurements to assist in creating a well-fitted device. Subsequently, the advent of plastics and later even more modern materials such as carbon fiber composites and aramid fibers as materials of choice for construction necessitated the idea of creating a plaster of Paris mold of the body part in question. This method is still extensively used throughout the industry. By introducing composite materials made of carbon fiber materials and aramid fibers embedded in an epoxy resin matrix, the weight of modern orthoses is extremely reduced. With this technique, modern orthoses can achieve perfect stiffness in the areas where this is necessary (e.g., the connection between the ankle and knee joint) and flexibility in the areas where flexibility is required (e.g., in the area of the forefoot on the foot part of an orthosis).
Semi-finished products – they are used for fast supply in the case of diseases that occur frequently. They are manufactured industrially and in some cases can be adapted to the anatomical body conditions. Semi-finished products are also referred to as prefabricated products and custom fitted products, but in these cases it is not custom-fabricated.
Finished products – these include short-term orthoses or bandages for a limited duration of therapy and are manufactured industrially. Finished products are also referred to as off-the-shelf products.
Both custom-fabricated products and semi-finished products are used in long-term care and are manufactured or adapted by the orthotist or by trained orthopedic technicians according to the prescription. In many countries the physician or clinician defines the functional deviations in his prescription, e.g. paralysis (paresis) of the calf muscles (M. Triceps Surae) and derives the indication from this, e.g. orthotic to restore safety when standing and walking after a stroke. The orthotist creates another detailed physical examination and compares it with the prescription from the physician. The orthotist describes the configuration of the orthosis, which shows which orthotic functions are required to compensate for the functional deviation of the neuromuscular or skeletal system and which functional elements must be integrated into the orthosis for this. Ideally, the necessary orthotic functions and the functional elements to be integrated are discussed in an interdisciplinary team between physician, physical therapist, orthotist and patient.
Lower limb orthoses
All orthoses that affect the foot, the ankle joint, the lower leg, the knee joint, the thigh or the hip joint belong to the category of orthoses for the lower extremities.
Paralysis orthoses
Paralysis orthoses are used for partial or complete paralysis, as well as complete functional failure of muscles or muscle groups, or incomplete paralysis (paresis). They are intended to correct or improve functional limitations or to replace functions that have been lost as a result of the paralysis. Functional leg length differences caused by paralysis can be compensated for by using orthosis.
For the quality and function of a paralysis orthosis, it is important that the orthotic shell is in total-contact with the patient's leg to create an optimal fit, which is why a custom-made orthotic is often preferred. As reducing the weight of an orthosis significantly lessens the energy needed to walk with it, the use of light weight and highly resilient materials such as carbon fiber, titanium and aluminum is indispensable for the manufacture of a custom-made orthosis.
The production of a custom-made orthotic also allows the integration of orthotic joints, which means the dynamics of the orthotic can be matched exactly with the pivot points of the patient's anatomical joints. As a result, the dynamics of the orthosis take place exactly where dictated by the patient's anatomy. Since the dynamics of the orthosis are executed via the orthotic joints, it is possible to manufacture the orthotic shells as stable and torsion-resistant, which is necessary for the quality and function of the orthosis. The orthosis thus offers the necessary stability to regain the security that has been lost due to paralysis when standing and walking.
In addition, an orthosis can be individually configured through the use of orthosis joints. In this way, the combination of the orthotic joints and the adjustability of the functional elements can be adjusted to compensate for any existing functional deviations that have resulted from the muscle weakness. The goal of a high-quality orthotic fitting is to adjust the functional elements so precisely that the orthosis provides the necessary support while restricting the dynamics of the lower extremities as little as possible to preserve the remaining functionality of the muscles.
Determination of strength levels for physical examination
In the case of paralysis due to disease or injury to the spinal/peripheral nervous system, a physical examination is needed to determine the strength levels of the affected leg's six major muscle groups and the orthosis's necessary functions.
The dorsiflexors move the foot through concentric muscle work around the axis of the ankle in the direction of dorsiflexion and control the plantar flexion through eccentric muscle work.
The plantar flexors contribute significantly to being able to stand upright by actuating the forefoot lever and thereby increasing the standing area when standing. This group of muscles moves the foot in the direction of plantar flexion.
The knee extensors extend the knee in the direction of the knee extension.
The knee flexors bend the knee in the direction of the knee flexion.
The hip flexors bend the hip joint toward the hip flexion.
The hip extensors stretch the hip joint in the direction of the hip extension and, at the same time, extend the knee in the direction of the knee extension.
According to Vladimir Janda, a muscle function test is carried out to determine strength levels. The degree of paralysis is given for each muscle group on a scale from 0 to 5, with the value 0 indicating complete paralysis (0%) and the value 5 indicating normal strength (100%). The values between 0 and 5 indicate a percentage reduction in muscle function. All strength levels below five are called muscle weakness.
The combination of strength levels of the muscle groups determines the type of orthosis (AFO or KAFO) and the functional elements necessary to compensate for restrictions caused by the reduced muscular strength levels.
Physical examination for paralysis due to diseases or injuries to the spinal cord and/or the peripheral nervous system
Paralysis may be caused by injury to the spinal or peripheral nervous system after spinal cord injury, or by diseases such as spina bifida, poliomyelitis and Charcot-Marie-Tooth disease. In these patients, knowledge of the strength levels of the large muscle groups is necessary to configure the orthotic for the necessary functions.
Physical examination for paralysis due to diseases or injuries to the central nervous system
Paralysis caused by diseases or injuries to the central nervous system (e.g. cerebral palsy, traumatic brain injury, stroke, and multiple sclerosis) can cause incorrect motor impulses that often result in clearly visible deviations in gait. The usefulness of muscle strength tests is therefore limited, as even with high degrees of strength, disturbances to the gait pattern can occur due to the incorrect control of the central nervous system.
Cerebral palsy and traumatic brain injury
In ambulatory patients with paralysis due to cerebral palsy or traumatic brain injury, the gait pattern is analysed as part of the physical examination in order to determine the necessary functions of an orthosis.
One way of classifying gait is according to the "Amsterdam Gait Classification", which describes five gait types. To assess the gait pattern, the patient is viewed directly, or via a video recording, from the side of the leg being assessed. At the point when the leg is mid-stance the knee angle and the contact of the foot with the ground are assessed. The five gait types are:
Type 1, the knee angle is normal and foot contact is complete.
Type 2, the knee angle is hyperextended and the foot contact is complete.
Type 3, the knee angle is hyperextended and foot contact is incomplete (only on the forefoot).
Type 4, the knee angle is flexed and foot contact is incomplete (only on the forefoot).
Type 5, the knee angle is flexed and foot contact is complete, this is also known as crouch gait.
Patients with paralysis due to cerebral palsy or traumatic brain injury are usually treated with an ankle-foot orthosis (AFO). Although in these patients the muscles are not paralyzed but being sent the wrong impulses from the brain, the functional elements used in the orthotics are the same for both groups. The compensatory gait is an unconscious reaction to the lack of security when standing or walking that usually worsens with increasing age; if the right functional elements are integrated into the orthosis to counter this, and maintain physiological mobility, the right motor impulses are sent to create new cerebral connections. The goal of an orthotic is the best possible approximation of the physiological gait pattern.
Stroke
In the case of paralysis after a stroke, rapid care with an orthosis is necessary. Often areas of the brain are affected that contain "programs" for controlling the musculoskeletal system. With the help of an orthosis, physiological standing and walking can be relearned, preventing long term health consequences caused by an abnormal gait pattern. According to Vladimir Janda, when configuring the orthotic it is important to understand that the muscle groups are not paralyzed, but are controlled by the brain with wrong impulses, and this is why a muscle function test can lead to incorrect results when assessing the ability to stand and walk.
An important basic requirement for regaining the ability to walk is that the patient trains early on to stand on both legs safely and well balanced. An orthosis with functional elements to support balance and safety when standing and walking can be integrated into physical therapy from the first standing exercises, and this makes the work of mobilizing the patient at an early stage easier. With the right functional elements that maintain physiological mobility and provide security when standing and walking, the necessary motor impulses to create new cerebral connections can occur. Clinical studies confirm the importance of orthoses in stroke rehabilitation.
Patients with paralysis after a stroke are often treated with an ankle-foot orthosis (AFO), as after a stroke stumbling can occur if only the dorsiflexors are supplied with incorrect impulses from the central nervous system. This can lead to insufficient foot lifting during swing phase of walking, and in these cases, an orthosis that only has functional elements to support the dorsiflexors can be helpful. Such an orthosis is also called drop foot orthosis. When configuring a foot lifter orthosis, adjustable functional elements for setting the resistance can be included, which make it possible to adapt the passive lowering of the forefoot (plantar flexion) to the eccentric work of the dorsal flexors during loading response.
In cases where the muscle group of the plantar flexors is supplied with wrong impulses from the central nervous system, which leads to uncertainty when standing and walking, an unconscious compensatory gait can occur. When configuring an orthosis functional elements that can restore safety when standing and walking must be used in these cases; a foot lifter orthosis is not suitable as it only compensates for the functional deviations caused by weakness of the dorsiflexors.
Patients with paralysis after stroke who are able to walk have the option of analysing the gait pattern in order to determine the optimal function of an orthosis. One way of assessing is the classification according to the "N.A.P. Gait Classification", which is a physiotherapeutic treatment concept. According to this classification, the gait pattern is assessed in the mid-stance phase and described as one of four possible gait types.
This assessment is a two step process; in the first step, the patient is viewed from the side of the leg to be assessed, either directly or via a video recording. In gait type 1 the knee angle is hyperextended, while in type 2, the knee angle is flexed. In the second step, the patient is viewed from the front to determine if the foot is inverted, if it is the letter "a" is added to the gait. This is associated with a varus deformity of the knee. If instead the patient stands on the inner edge of the foot (eversion), which is associated with a valgus deformity of the knee, the letter "b" is added to the gait type. Patients are thus classified as gait types 1a, 1b, 2a or 2b. The goal of orthotic fitting for patients who are able to walk is the best possible approximation of the physiological gait pattern.
Multiple sclerosis (MS)
In the case of paralysis due to multiple sclerosis, the degree of strength of the six major muscle groups of the affected leg should be determined as part of the physical examination in order to determine the necessary functions of an orthosis, just as in the case of diseases or injuries to the spinal/peripheral nervous system. However, patients with multiple sclerosis may experience muscular fatigue as well. The fatigue can be more or less pronounced and, depending on the severity, can lead to considerable restrictions in everyday life. Persistent stress, such as from walking, causes a deterioration in muscle function and has a significant effect on the spatial and temporal parameters of walking, for example by significantly reducing the cadence and walking speed. Fatigue can be measured as muscle weakness. When determining the strength levels of the six major muscle groups as part of the patient's medical history, fatigue can be taken into account by using a standardized six-minute walking test. According to Vladimir Janda the muscle function test is carried out in combination with the six-minute walk test in the following steps:
First muscle function test (without muscular fatigue)
Six-minute walk test directly followed by
Second muscle function test (with muscular fatigue)
This sequence of muscle function test and six-minute walk test is used to determine whether muscular fatigue can be induced. If the test reveals muscular fatigue, the strength levels and measured fatigue should be included in the planning of an orthosis, and when determining the functional elements.
Functional deviations in the case of paralysis of large muscle groups
Paralysis of the dorsiflexors – weakness of the dorsiflexors results in a drop foot. The patient's foot cannot be sufficiently lifted during the swing phase while walking, as the necessary concentric work of the dorsiflexors can not be activated. There is a risk of stumbling, and the patient cannot influence the shock absorption when walking (gait phase, loading response), as the eccentric work of the dorsiflexors is limited. After initial heel contact the forefoot either slaps too quickly on the floor via the heel rocker, which creates an audible noise, or the foot does touch the floor with forefoot first, which disrupts gait development.
Paralysis of the plantar flexors – If the plantar flexors are weak, the muscles of the forefoot lever are either inadequately activated or not activated at all. The patient has no balance when standing and has to support themself with aids such as crutches. The forefoot lever required for energy-saving walking in the gait phases from mid-stance to pre-swing cannot be activated by the plantar flexors. This leads to excessive dorsiflexion in the ankle joint in terminal stance and a loss of energy while walking. The center of gravity of the body lowers towards the end of the stance phase and the knee of the contralateral leg is flexed excessively. With each step, the center of gravity must be raised above the leg by straightening the excessively flexed knee. Since the plantar flexors originate above the knee joint, they also have a knee-extension effect in the stance phase.
Paralysis of the knee extensors – if the knee extensors are weak, there is an increased risk of falling when walking, as between loading response to the mid-stance the knee extensors control knee flexion inadequately, or not at all. To control the knee, the patient develops compensatory mechanisms that lead to an incorrect gait pattern, for example by exaggerated activation of the plantar flexors, leading into hyperextension of the knee, or when initial contact is with the forefoot and not the heel in order to prevent the knee-flexing effect of the heel rocker.
Paralysis of the knee flexors – if the knee flexors are weak, it is more difficult to flex the knee in pre-swing.
Paralysis of the hip flexors – if the hip flexors are weak, it is more difficult to flex the knee in pre-swing.
Paralysis of the hip extensors – the hip extensors help control of the knee against unwanted flexion when walking between loading response and mid-stance.
Functional elements in paralysis of large muscle groups
The functional elements of an orthosis ensure the flexion and extension movements of the ankle, knee and hip joints. They correct and control the movements and secure the joints against undesired incorrect movements, and help avoid falls when standing or walking.
Functional elements in paralysis of the dorsiflexors – if the dorsiflexors are weak, an orthosis should lift the forefoot during the swing phase in order to reduce the risk of the patient stumbling. An orthosis that has only one functional element for lifting the forefoot in order to compensate for a weakness in the dorsiflexors is also known as a drop foot orthosis. An AFO of the drop foot orthosis type is therefore not suitable for the care of patients with weakness in other muscle groups, as these patients require additional functional elements to be taken into account. Initial contact with the heel should be achieved by lifting the foot through the orthosis, and if the dorsiflexors are very weak, control of the rapid drop of the forefoot should be taken over by dynamic functional elements that allow for adjustable resistance of plantar flexion. Orthoses should be adapted to the functional deviation of the dorsiflexors in order to correct the shock absorption of the heel rocker lever during loading response, but should not block plantar flexion of the ankle joint as this leads to excessive flexion in the knee and hip and an increase in the energy needed for walking. This is why static functional elements are not recommended when there are newer technical alternatives.
Functional elements in paralysis of the plantar flexors – in order to compensate for a weakness of the plantar flexors, the orthosis has to transfer large forces that the strong muscle group would otherwise take over. These forces are transmitted in a similar way to a ski boot during downhill skiing via the functional elements of the foot part, ankle joint and lower leg shell. Dynamic functional elements are preferable for the ankle joint as static functional elements would completely block the dorsiflexion, which would have to be compensated for by the upper body, resulting in an increased energy cost when walking. The functional element's resistance to protect against unwanted dorsiflexion should be able to be adapted according to the weakness of the plantar flexors. In the case of very weak plantar flexors, the functional element's resistance against undesired dorsiflexion must be very high in order to compensate for the functional deviations this causes. Adjustable functional elements allow the resistance to be adjusted exactly to the weakness of the muscle, and scientific studies recommend adjustable resistance in patients with paralysis or weakness of the plantar flexors.
Functional elements in paralysis of knee extensors and hip extensors – in the case of weak knee extensors or hip extensors, the orthosis must take over the stability and stance phase control when walking. Different knee-securing functional elements are needed depending on the weakness of these muscles. In order to compensate for functional deviations with slightly weakness of these muscle groups, a free moving mechanical knee joint with the mechanical pivot point behind the anatomical knee pivot point can be sufficient. In the case of significant weakness, knee flexion when walking must be controlled by functional elements that mechanically secure the knee joint during the early stance phases between loading response and mid stance. Stance phase control knee joints which lock the knee in the early stance phases and release it for knee flexion during the swing phase can be used here, with these joints, a natural gait pattern can be achieved despite mechanically securing against unwanted knee flexion. In these cases, locked knee joints are often used, and while they have a good safety function, the knee joint remains mechanically locked during the swing phase while walking. Patients with locked knee joints have to manage the swing phase with a stiff leg, which only works if the patient develops compensatory mechanisms, such as by raising the body's center of gravity in the swing phase (Duchenne limping) or by swinging the orthotic leg to the side (circumduction). Stance phase control knee joints and locked joints can both be mechanically "unlocked" so the knee can be flexed to sit down.
Ankle–foot orthoses (AFO) in the field of paralysis orthoses
AFO is the abbreviation for ankle-foot orthoses, which is the English name for an orthosis that spans the ankle and foot. In the treatment of paralyzed patients, they are mainly used when there is a weakness of the dorsiflexors or plantar flexors.
Through the use of modern materials, such as carbon fibers and aramid fibers, and new knowledge about processing these materials into composite materials, the weight of orthotics has been reduced significantly. In addition to the weight reduction, these materials and technologies have created the possibility of making some areas of an orthosis so rigid that it can take over the forces of the weakened muscles (e.g. the connection from the ankle joint to the frontal contact surface on the shin), while at the same time leaving areas requiring less support very flexible (e.g. the flexible part of the forefoot).
It is now possible to combine the required rigidity of the orthotic shells with the dynamics in the ankle, with this, other new technologies, and the possibility of producing lightweight but rigid orthoses, new demands have been made of orthotics:
Despite the necessary rigidity, the orthoses should not block the mobility of the ankle.
Despite the necessary rigidity, the orthoses should not block the functionality of the muscles, but rather promote it.
Despite the necessary rigidity, contractures and spasticity should not be stimulated.
A custom-made AFO can compensate for functional deviations of muscle groups, it should be configured according to the patient data through a function and load calculation so that it meets the functional and load requirements. In calculating or configuring an AFO, variants are optimally matched to individual requirements for the functional elements of the ankle joint, for the stiffness of the foot shell, and for the shape of the lower leg shell. The size of these components is selected by matching their resilience to the load data.
An ankle joint based on new technology is the connection between the foot shell and the lower leg shell and at the same time contains all the necessary adjustable functional elements of an AFO.
Depending on the combination of the degree of paralysis of the dorsiflexors or plantar flexors, different functional elements to compensate for their weakness can be integrated into the ankle joint; if both muscle groups are affected, the elements should be integrated into one orthotic joint. The necessary dynamics and resistance to movements in the ankle can be adapted via adjustable functional elements in the ankle joint of the orthosis, which allows it to compensate for muscle weaknesses, provide safety when standing and walking, and still allow as much mobility as possible. For example, adjustable spring units with pre-compression can enable an exact adaptation of both static and dynamic resistance to the measured degree of muscle weakness. Studies show the positive effects of these new technologies. It is of great advantage if the resistances for these two functional elements can be set separately.
An AFO with functional elements to compensate for a weakness of the plantar flexors can also be used for slight weakness of the knee-securing muscle groups, the knee extensors and the hip extensors.
A drop foot orthosis is an AFO that only has one functional element for lifting the forefoot in order to compensate for a weakness in the dorsiflexors. If other muscle groups, such as the plantar flexors, are weak, additional functional elements must be taken into account, making a drop foot orthosis unsuitable for patients with weakness in other muscle groups.
In 2006, before these new technologies were available, the International Committee of the Red Cross published in its 2006 Manufacturing Guidelines for Ankle-Foot Orthoses, with the aim of providing people with disabilities worldwide standardized processes for the production of high-quality, modern, durable and economical devices.
Because new technologies are not widely used, AFOs are often made from polypropylene-based plastic, mostly in the shape of a continuous "L" shape, with the upright part behind the calf and the lower part under the foot, however, this only offers the rigidity of the material. AFOs made of polypropylene are still called "DAFO" (dynamic ankle-foot orthosis), "SAFO" (solid ankle-foot orthosis) or "Hinged AFO". DAFOs are not stable enough to transfer the high forces required to balance the weak plantar flexors when standing and walking, and SAFOs block the mobility of the ankle joint. A "Hinged AFO" only allowed for the compensation that could be achieved with the orthotic joints of the time, for example, they commonly block plantar flexion, as the joints cannot simultaneously transmit the large forces that are required to compensate for muscle deviations while also offering the necessary dynamics.
While there was a multitude of AFOs with differing designs in clinical practice, there was also a clear lack of details regarding the design and the materials used for manufacture, leading Eddison and Chockalingam to call for a new standardization of the terminology. With a focus on caring for children with cerebral palsy there is a recommendation to investigate the potential for gait pattern improvement via the design and manufacture of orthotics made of polypropylene. On the other hand, integrating orthotic joints with modern functional elements into the production of older technologies using polypropylene is unusual because the orthotic shells made of polypropylene either could not transfer the high forces or would be too soft.
New studies now show the better possibilities for improving the gait pattern through the new technologies.
The International Committee of the Red Cross published its manufacturing guidelines for ankle–foot orthoses in 2006, and, unfortunately, today's terminologies are still based those guidelines and therefore require a particularly high level of explanation. The intent was to provide standardized procedures for the manufacture of high-quality modern, durable and economical devices to people with disabilities throughout the world. However, with the new technologies available, the main types mentioned are in need of revision today.
Knee-ankle-foot orthosis (KAFO) in the field of paralysis orthoses
KAFO is the abbreviation for knee-ankle-foot orthoses, which spans the knee, ankle and foot. In the treatment of paralyzed patients, a KAFO is used when there is a weakness of the knee or hip extensors. They have two orthotic joints: an ankle joint between the foot and lower leg shells and a knee joint between the lower leg and thigh shells.
KAFOs can be roughly divided into three variants, depending on whether the mechanical knee joint is: locked, unlocked or locked and unlocked.
KAFO with locked knee joint - The mechanical knee joint is locked both when standing and also when walking (in both the stance and swing phases) in order to achieve the necessary stability. To sit, the user can unlock the knee joint. When walking with a locked knee joint it is difficult for the user to swing the leg forward and, in order to not stumble, the leg must be swung forward and out in a circular arc (circumduction) or the hip must be raised unnaturally to swing the stiff leg. Each of these incorrect gait patterns can lead to secondary diseases in the bone and muscle system, and such compensatory movement patterns lead to increased energy consumption when walking. The film Forrest Gump impressively shows how the main character Forrest Gump is additionally hindered in his urge to move by such orthoses. For centuries, KAFOs were built with mechanical knee joints that stiffened the knee of the paralyzed leg, and even today, such orthotic fittings are still common. Typical designations for a KAFO with a locked knee joint include "KAFO with Swiss lock" or "KAFO with drop lock lock".
KAFO with unlocked knee joint - An unlocked knee joint can move freely both when standing and when walking, both in the stance phase and in the swing phase. In order for the leg to swing through without stumbling, knee flexion of approximately 60° is allowed; the user does not need to unlock the knee joint to sit. As a KAFO with an unlocked knee joint can provide only minor compensation for paralysis-related issues while standing and walking, an orthotic knee joint with a rearward displacement of the pivot point can be installed in order to increase safety. However, even with this, a KAFO with a non-locked knee joint should only be used in cases of minor paralysis of the knee and hip extensors. With more severe paralysis and low levels of strength in these muscle groups, there is a significant risk of falling. A typical designation for a KAFO with a unlocked knee joint is, among other things, "KAFO with knee joint for movement control".
KAFO with locked and unlocked knee joint - The mechanical knee joint of a KAFO with locked and unlocked knee joint is locked when walking in the stance phase, providing the necessary stability and security for the user. The knee joint is then automatically unlocked in the swing phase, allowing the leg to be swung through without stumbling. In order to be able to walk efficiently, without stumbling, and without compensating mechanisms, the joint should allow knee flexion of approximately 60° in the swing phase. The first promising developments of automatic knee joints, or stance phase locking knee joints, emerged in the 1990s. In the beginning there were automatic mechanical constructions that took over the locking and unlocking, now automatic electromechanical and automatic electrohydraulic systems are available that make standing and walking safer and more comfortable. Various terms are used for a KAFO with a locked and unlocked knee joint. Typical designations are "KAFO with automatic knee joint" or "KAFO with stance phase control knee joint". In scientific articles, the English term Stance Control Orthoses SCO is often used, but as this term differs from the ICS classification, one of the first two terms is preferable.
Different functional elements to compensate for weakness of the dorsiflexors or plantar flexors can be integrated into the ankle joint of the orthosis depending on the degree of paralysis of the two muscle groups. It is of great advantage if the resistances for these two functional elements can be set separately. The functional elements to compensate for paralysis of the knee-securing muscle groups of the knee and hip extensors are integrated into the knee joint of the orthosis via knee-securing functional elements. A KAFO can use a variety of combinations of different variants in the stiffness of the foot shell, the different variants of the functional elements of a dynamic ankle joint, the variants in the shape of the lower leg shell, and the functional elements of a knee joint to compensate for the user's limitations.
Hip-knee-ankle-foot orthosis (HKAFO) in the field of paralysis orthoses
HKAFO is the abbreviation for hip-knee-ankle-foot orthoses; which is the English name for an orthosis that spans the hip, the knee, the ankle and the foot. In the treatment of paralyzed patients, a HKAFO is used when there is a weakness of the pelvic stabilizing trunk muscles.
Relief orthoses
Relief Orthoses are used when there is degeneration to a joint (from "wear and tear" for example) or after an injury such as a torn ligament. Relief orthoses are also used after operations such as operations on the joint ligaments, other bony, muscular structures, or after a complete replacement of a joint.
Relief orthosis may also be used to:
Control, guide, limit and/or immobilize an extremity, joint or body segment for a particular reason
Restrict movement in a given direction
Assist movement generally
Reduce weight-bearing forces for a particular purpose
Aid rehabilitation from fractures after the removal of a cast
Otherwise correct the shape and/or function of the body, to provide easier movement capability or reduce pain
Ulcer healing orthoses (UHO)
A custom-made ankle/foot orthosis can be used for the treatment of patients with foot ulcers, it is a rigid L-shaped support member with a rigid anterior support shell on an articulated hinge. The plantar portion of the L-shaped member has at least one ulcer-protecting hollow to allow the user to transfer their weight away from the ulcer to facilitate treatment. The anterior support shell is designed with a lateral hinged attachment to take advantage of the medial tibial flare structure to enhance the weight-bearing properties of the orthosis. A flexible, polyethylene hinge attaches the support shell to the L-shaped member and straps securely attach the anterior support shell to the user's lower leg.
Foot orthoses (FO)
Foot orthoses (commonly called orthotics) are devices inserted into shoes to provide support for the foot by redistributing ground reaction forces acting on the foot joints while standing, walking or running. They may be either pre-moulded (also called pre-fabricated) or custom made according to a cast or impression of the foot. They are used by everyone from athletes to the elderly to accommodate biomechanical deformities and a variety of soft tissue conditions. Foot orthoses are effective at reducing pain for people with painful high-arched feet, and may be effective for people with rheumatoid arthritis, plantar fasciitis, first metatarsophalangeal (MTP) joint pain or hallux valgus (bunions). For children with juvenile idiopathic arthritis (JIA) custom-made and pre-fabricated foot orthoses may also reduce foot pain. Foot orthoses may also be used in conjunction with properly fitted orthopedic footwear in the prevention of diabetic foot ulcers. A real-time weight bearing orthotic can be created using a neutral position casting device and the Vertical Foot Alignment System VFAS.
Ankle–foot orthoses (AFO) in the field of relief orthoses
An AFO can also be used to immobilize the ankle and lower leg in the presence of arthritis or a fracture. Ankle–foot orthoses are the most commonly used orthoses, making up about 26% of all orthoses provided in the United States. According to a review of Medicare payments from 2001 to 2006, the base cost of an AFO was about $500 to $700.
Knee orthoses (KO) in the field of relief orthoses
A knee orthosis (KO) or knee brace extends above and below the knee joint and is generally worn to support or align the knee. In the case of diseases causing neurological or muscular impairment of muscles surrounding the knee, a KO can prevent flexion, extension, or instability of the knee. If the ligaments or cartilage of the knee are affected, a KO can provide stabilization to the knee by replacing their functions. For instance, knee braces can be used to relieve pressure from diseases such as arthritis or osteoarthritis by realigning the knee joint. In this way a KO may help reduce osteoarthritis pain, however, there is no clear evidence about the most effective orthosis or the best approach to rehabilitation. A knee brace is not meant to treat an injury or disease on its own, but is used as a component of treatment along with drugs, physical therapy and possibly surgery. When used properly, a knee brace may help an individual to stay active by enhancing the position and movement of the knee or reducing pain.
Prophylactic, functional and rehabilitation braces
Prophylactic braces are used primarily by athletes participating in contact sports. Evidence indicates that prophylactic knee braces, like the ones football linemen wear that are often rigid with a knee hinge, are ineffective in reducing anterior cruciate ligament tears, but may be helpful in resisting medial and lateral collateral ligament tears.
Functional braces are designed for use by people who have already experienced a knee injury and need support while recovering from it, or to help people who have pain associated with arthritis. They are intended to reduce the rotation of the knee, support stability, reduce the chance of hyperextension, and increase the agility and strength of the knee. The majority of these are made of elastic. They are the least expensive of all braces and are easily found in a variety of sizes.
Rehabilitation braces are used to limit the movement of the knee in both medial and lateral directions, these braces often have an adjustable range of motion, and can be used to limit flexion and extension following ACL reconstruction. They are primarily used after injury or surgery to immobilize the leg and are larger in size than other braces, due to their function.
Soft braces
A soft brace, sometimes called soft support or a bandage, belong to the field of orthoses and are supposed to protect the joints from excessive loads.
Soft braces are also classified according to regions of the body. In sport, bandages are used to protect bones and joints, and prevent and protect injuries. Bandages should also allow proprioception. They mostly consist of textiles, some of which have supportive elements. The supporting functions are low compared to paralysis and relief orthoses, though they are sometimes used prophylactically or to optimize performance in sport. At present, the scientific literature does not provide sufficient high quality research to allow for strong conclusions on their effectiveness and cost-effectiveness.
Upper limb orthoses
Upper-limb (or upper extremity) orthoses are mechanical or electromechanical devices applied externally to the arm, or segments of it, in order to restore or improve function or structural characteristics of the arm segments enclosed in the device. In general, musculoskeletal problems that may be alleviated by the use of upper limb orthoses include those resulting from trauma or disease (arthritis for example). They may also benefit individuals who have a neurological impairment from a stroke, spinal cord injury, or peripheral neuropathy.
Types of upper-limb orthoses
Upper-limb orthoses
Clavicular and shoulder orthoses
Arm orthoses
Functional arm orthoses
Elbow orthoses
Forearm-wrist orthoses
Forearm-wrist-thumb orthoses
Forearm-wrist-hand orthoses
Hand orthoses
Upper-extremity orthoses (with special functions)
Spinal orthoses
Scoliosis, a condition describing an abnormal curvature of the spine, may in certain cases be treated with spinal orthoses, such as the Milwaukee brace, Boston brace, Charleston bending brace, or Providence brace. As scoliosis most commonly develops in adolescent females who are undergoing their adolescent growth spurt, compliance is hampered by patient concerns about appearance and movement restrictions caused by the brace.
Spinal orthoses may also be used in the treatment of spinal fractures. A Jewett brace, for instance, may be used to aid healing of an anterior wedge fracture involving the T10 to L3 vertebrae, and a body jacket may be used to stabilize more involved fractures of the spine. There are several types of orthoses for managing cervical spine pathology. The halo brace is the most restrictive cervical thoracic orthosis in use; it is used to immobilize the cervical spine, usually following fracture, and was developed by Vernon L. Nickel at Rancho Los Amigos National Rehabilitation Center in 1955.
Head orthoses
Helmets such as a cranial molding orthoses is an example of orthoses for the head. These devices are often suggested for infants with positional plagiocephaly.
| Biology and health sciences | Fields of medicine | Health |
22395556 | https://en.wikipedia.org/wiki/Agricultural%20engineering | Agricultural engineering | Agricultural engineering, also known as agricultural and biosystems engineering, is the field of study and application of engineering science and designs principles for agriculture purposes, combining the various disciplines of mechanical, civil, electrical, food science, environmental, software, and chemical engineering to improve the efficiency of farms and agribusiness enterprises as well as to ensure sustainability of natural and renewable resources.
An agricultural engineer is an engineer with an agriculture background. Agricultural engineers make the engineering designs and plans in an agricultural project, usually in partnership with an agriculturist who is more proficient in farming and agricultural science.
History
The first use of agricultural engineering was the introduction of irrigation in large scale agriculture in the Nile and the Euphrates rivers before 2000 B.C. Large irrigation structures were also present in Baluchistan and India before Christian era. In other parts of Asia, agricultural engineering was heavily present in China. In South America irrigation was practiced in Peru by the Incas and in North America by the Aztecs.
The earliest plough was the ard or scratch-plough.
Settlers practiced irrigation in the vicinity of San Antonio in 1715, the Mormons practiced irrigation in Salt Lake Valley in 1847.
With growing mechanization and steam power in the Industrial Revolution, a new age in agricultural engineering began. Over the course of the Industrial Revolution, mechanical harvesters and planters would replace field hands in most of the food and cash crop industries. Mechanical threshing was introduced in 1761 by John Lloyd, Magnus Strindberg and Dietrich. Beater bar threshing machine was built by Andrew Meikle in 1786. A cast iron plow was first made by Charles Newbold between 1790 and 1796.
James Smith constructed a mower in 1811. George Berry used a steam combine harvester in 1886. John Deere made his first steel plow in 1833. The two horse cultivator was first about 1861.
The introduction of these engineering concepts into the field of agriculture allowed for an enormous boost in the productivity of crops, dubbed a "second agricultural revolution" which consisted of:
Shift from peasant subsistence-farming to cash-farming for the market
Technical changes of crop rotations and livestock improvement
Labour being replaced by machinery
In the 20th century, with the rise in reliable engines in airplanes, cropdusters were implemented to disperse pesticides. Benjamin Holt built a combine harvester powered by petrol in 1911. Erwin Peucker constructed bulldog tractors 1936. Deutz-Fahr produced the rotary hay tedder in 1961.
In the late 20th century, genetically modified foods (GMOs) were created, giving another large boost to crop yields and resistance to pests.
Sub-disciplines
Agricultural engineering has many sub-disciplines, the most common of which are listed here:
Agricultural machinery
Agricultural structures
Agricultural surveying
Aquaculture
Biomechanics and ergonomics
Forestry engineering
Irrigation
Land development
Pesticides
Precision agriculture
Soil management
Roles of agricultural engineers
Agricultural engineers may perform tasks such as planning, supervising, and managing the building of dairy effluent schemes, irrigation, drainage, flood water control systems, performing environmental impact assessments, agricultural product processing and interpret research results and implement relevant practices. A large percentage of agricultural engineers work in academia or for government agencies. Some are consultants employed by private engineering firms, while others work in industry for manufacturers of agricultural machinery, equipment, processing technology, and structures for housing livestock and storing crops. Agricultural engineers work in production, sales, management, research and development, or applied science.
Armenia
In 2006, Armenia's agricultural sector accounted for about 20 percent of the GDP. By 2010, it grew to about 25 percent. This was and is higher than in Armenia's neighboring countries of Georgia, Azerbaijan, Turkey, and Iran, in which the contribution of agriculture to the GDP in 2017 was 6.88, 5.63, 6.08, and 9.05 percent, respectively.
Philippines
In the Philippines, the professional designation is registered agricultural and biosystems engineer'''. These engineers are licensed and accredited after successfully passing the Agricultural and Biosystems Engineering Licensure Examination. A prospective agricultural and biosystems engineer must have a four-year Bachelor of Science in Agricultural and Biosystems Engineering.
The practice of agricultural and biosystems engineering also includes the following:
Consultation, valuation, investigation, and management services on agricultural and biosystems engineering;
Management or supervision and the preparation of engineering designs, plans, specifications, project studies, and estimates for agricultural and biosystems, aquaculture and fishery, forest product machinery, agricultural and biosystems buildings and structures, farm electrification and energy systems, agricultural and biosystems processing equipment, irrigation, and soils conservation systems and facilities, agricultural and biosystems waste utilization systems and facilities;
Conducting research and development, training and extension work, and consultancy services on agricultural and biosystems engineering facilities/services, systems, and technologies;
Testing, evaluation, and inspection of agricultural and biosystems, fishery and forest product machinery, and other related agricultural and biosystems engineering facilities and equipment.
Management, manufacturing, and marketing of agricultural and biosystems machinery and other related agricultural and biosystems engineering facilities and equipment;
Teaching agricultural and biosystems engineering subjects in institutions of learning in the Philippines;
Employment with the government provided such item or position requires the knowledge and expertise of an agricultural and biosystems engineer.
United Kingdom
In the United Kingdom, the term agricultural engineer is often also used to describe a person who repairs or modifies agricultural equipment.
United States
The American Society of Agricultural Engineers, now known as the American Society of Agricultural and Biological Engineers (ASABE), was founded in 1907. It is a leading organization in the agricultural engineering field. The ASABE provides safety and regulatory standards for the agricultural industry. These standards and regulations are developed on an international scale for fertilizers, soil conditions, fisheries, biofuels, biogas, feed machinery, tractors, and machinery.
| Technology | Disciplines | null |
6303947 | https://en.wikipedia.org/wiki/Banana%20pepper | Banana pepper | The banana pepper (also known as the yellow wax pepper or banana chili) is an average-sized member of the chili pepper family that has a mild, tangy taste. While typically bright yellow, it is possible for them to change to green, red, or orange as they ripen. It is often pickled, stuffed or used as a raw ingredient in foods. It is a cultivar of the species Capsicum annuum. Its flavor is not very hot (0–500 Scoville units) and, as is the case with most peppers, its heat depends on the maturity of the pepper, with the ripest being sweeter than younger ones.
Nomenclature
A mature fruit will be about 2–3 inches (5–8 cm) in length and have a curved shape and yellowish colour similar to a banana, giving rise to the fruit's common name. Friggitelli (pepperoncini) are often erroneously referred to as banana peppers. The hot varieties of banana pepper are called Hungarian wax peppers.
Cultivation
The plant requires full sun, like other Capsicum annuum varieties, and should be treated the same as most other plants in the pepper family. Plants can be grown from seed and cuttings. A mature plant will reach 1 to 2 feet tall and can be grown in many climates, but prefers warmer climates. Cultivars include Early Sweet Banana, Hungarian Yellow Wax, Long Sweet Yellow, Sweet Banana, and Sweet Hungarian.
Nutritional information
Raw banana peppers contain 92% water, 5% carbohydrates, and negligible fat and protein (table). They are rich sources of vitamin C, containing 92% of the Daily Value (DV) in a 100 gram reference amount (table). Vitamin B6 is present in substantial content of 21% DV, with no other micronutrients in appreciable amounts (table).
Uses
Culinary
Pickled banana peppers are commonly sold sliced and used to garnish pizzas, sandwiches and Greek salads.
Pickled and stuffed banana peppers are common inclusions on antipasto bars filled with prosciutto and/or cheese.
Stuffed banana peppers are served warm with a variety of Italian sausage and cheeses.
Chopped or diced banana peppers are used in many relishes and salsas to add sweetness with other peppers providing heat.
Banana peppers may be jellied along with other hot green peppers such as jalapeños.
Bhaji, a snack food.
| Biology and health sciences | Botanical fruits used as culinary vegetables | Plants |
6308405 | https://en.wikipedia.org/wiki/Delta%20potential | Delta potential | In quantum mechanics the delta potential is a potential well mathematically described by the Dirac delta function - a generalized function. Qualitatively, it corresponds to a potential which is zero everywhere, except at a single point, where it takes an infinite value. This can be used to simulate situations where a particle is free to move in two regions of space with a barrier between the two regions. For example, an electron can move almost freely in a conducting material, but if two conducting surfaces are put close together, the interface between them acts as a barrier for the electron that can be approximated by a delta potential.
The delta potential well is a limiting case of the finite potential well, which is obtained if one maintains the product of the width of the well and the potential constant while decreasing the well's width and increasing the potential.
This article, for simplicity, only considers a one-dimensional potential well, but analysis could be expanded to more dimensions.
Single delta potential
The time-independent Schrödinger equation for the wave function of a particle in one dimension in a potential is
where is the reduced Planck constant, and is the energy of the particle.
The delta potential is the potential
where is the Dirac delta function.
It is called a delta potential well if is negative, and a delta potential barrier if is positive. The delta has been defined to occur at the origin for simplicity; a shift in the delta function's argument does not change any of the following results.
Solving the Schrödinger equation
Source:
The potential splits the space in two parts ( and ). In each of these parts the potential is zero, and the Schrödinger equation reduces to
this is a linear differential equation with constant coefficients, whose solutions are linear combinations of and , where the wave number is related to the energy by
In general, due to the presence of the delta potential in the origin, the coefficients of the solution need not be the same in both half-spaces:
where, in the case of positive energies (real ), represents a wave traveling to the right, and one traveling to the left.
One obtains a relation between the coefficients by imposing that the wavefunction be continuous at the origin:
A second relation can be found by studying the derivative of the wavefunction. Normally, we could also impose differentiability at the origin, but this is not possible because of the delta potential. However, if we integrate the Schrödinger equation around , over an interval :
In the limit as , the right-hand side of this equation vanishes; the left-hand side becomes
because
Substituting the definition of into this expression yields
The boundary conditions thus give the following restrictions on the coefficients
Bound state (E < 0)
In any one-dimensional attractive potential there will be a bound state. To find its energy, note that for , is imaginary, and the wave functions which were oscillating for positive energies in the calculation above are now exponentially increasing or decreasing functions of x (see above). Requiring that the wave functions do not diverge at infinity eliminates half of the terms: . The wave function is then
From the boundary conditions and normalization conditions, it follows that
from which it follows that must be negative, that is, the bound state only exists for the well, and not for the barrier. The Fourier transform of this wave function is a Lorentzian function.
The energy of the bound state is then
Scattering (E > 0)
For positive energies, the particle is free to move in either half-space: or . It may be scattered at the delta-function potential.
The quantum case can be studied in the following situation: a particle incident on the barrier from the left side . It may be reflected or transmitted .
To find the amplitudes for reflection and transmission for incidence from the left, we put in the above equations (incoming particle), (reflection), (no incoming particle from the right) and (transmission), and solve for and even though we do not have any equations in .
The result is
Due to the mirror symmetry of the model, the amplitudes for incidence from the right are the same as those from the left. The result is that there is a non-zero probability
for the particle to be reflected. This does not depend on the sign of , that is, a barrier has the same probability of reflecting the particle as a well. This is a significant difference from classical mechanics, where the reflection probability would be 1 for the barrier (the particle simply bounces back), and 0 for the well (the particle passes through the well undisturbed).
The probability for transmission is
Remarks and application
The calculation presented above may at first seem unrealistic and hardly useful. However, it has proved to be a suitable model for a variety of real-life systems.
One such example regards the interfaces between two conducting materials. In the bulk of the materials, the motion of the electrons is quasi-free and can be described by the kinetic term in the above Hamiltonian with an effective mass . Often, the surfaces of such materials are covered with oxide layers or are not ideal for other reasons. This thin, non-conducting layer may then be modeled by a local delta-function potential as above. Electrons may then tunnel from one material to the other giving rise to a current.
The operation of a scanning tunneling microscope (STM) relies on this tunneling effect. In that case, the barrier is due to the air between the tip of the STM and the underlying object. The strength of the barrier is related to the separation being stronger the further apart the two are. For a more general model of this situation, see Finite potential barrier (QM). The delta function potential barrier is the limiting case of the model considered there for very high and narrow barriers.
The above model is one-dimensional while the space around us is three-dimensional. So, in fact, one should solve the Schrödinger equation in three dimensions. On the other hand, many systems only change along one coordinate direction and are translationally invariant along the others. The Schrödinger equation may then be reduced to the case considered here by an Ansatz for the wave function of the type .
Alternatively, it is possible to generalize the delta function to exist on the surface of some domain D (see Laplacian of the indicator).
The delta function model is actually a one-dimensional version of the Hydrogen atom according to the dimensional scaling method developed by the group of Dudley R. Herschbach
The delta function model becomes particularly useful with the double-well Dirac Delta function model which represents a one-dimensional version of the Hydrogen molecule ion, as shown in the following section.
Double delta potential
The double-well Dirac delta function models a diatomic hydrogen molecule by the corresponding Schrödinger equation:
where the potential is now
where is the "internuclear" distance with Dirac delta-function (negative) peaks located at (shown in brown in the diagram). Keeping in mind the relationship of this model with its three-dimensional molecular counterpart, we use atomic units and set . Here is a formally adjustable parameter. From the single-well case, we can infer the "ansatz" for the solution to be
Matching of the wavefunction at the Dirac delta-function peaks yields the determinant
Thus, is found to be governed by the pseudo-quadratic equation
which has two solutions . For the case of equal charges (symmetric homonuclear case), , and the pseudo-quadratic reduces to
The "+" case corresponds to a wave function symmetric about the midpoint (shown in red in the diagram), where , and is called gerade. Correspondingly, the "−" case is the wave function that is anti-symmetric about the midpoint, where , and is called ungerade (shown in green in the diagram). They represent an approximation of the two lowest discrete energy states of the three-dimensional H2^+ and are useful in its analysis. Analytical solutions for the energy eigenvalues for the case of symmetric charges are given by
where W is the standard Lambert W function. Note that the lowest energy corresponds to the symmetric solution . In the case of unequal charges, and for that matter the three-dimensional molecular problem, the solutions are given by a generalization of the Lambert W function (see ).
One of the most interesting cases is when qR ≤ 1, which results in . Thus, one has a non-trivial bound state solution with . For these specific parameters, there are many interesting properties that occur, one of which is the unusual effect that the transmission coefficient is unity at zero energy.
| Physical sciences | Quantum mechanics | Physics |
6309425 | https://en.wikipedia.org/wiki/Tamandua | Tamandua | Tamandua is a genus of anteaters in the Myrmecophagidae family with two species: the southern tamandua (T. tetradactyla) and the northern tamandua (T. mexicana). They live in forests and grasslands, are semiarboreal, and possess partially prehensile tails. They mainly eat ants and termites, but they occasionally eat bees, beetles, and insect larvae. In captivity, they will eat fruits and meat. They have no teeth and depend on their powerful gizzards to break down their food.
Taxonomy
The genus name is derived from the word tamandûá in Tupi first recorded by Joseph of Anchieta in his Epistola quam plurimarum rerum naturalium quae S. Vicentii (nunc S. Pauli) provinciam incolunt sistens descriptionem (1560); ta- in tamãdu'á is deduced from taixi "ant" and the other half from mondé "to catch", mondá "thief" or monduár "hunter".
Extant species
Description
Tamanduas have tapered heads with a long, tubular snout, small eyes, and protruding ears. Their tapered mouths house a tongue reaching upwards of in length. The tail is hairless and pink in color, marked with an irregular pattern of black blotches. The forefeet possess four clawed digits, the third digit bearing the largest claw, while the hind feet have five digits.
Their fur is thick, bristly, yellowish-white to fawn in color, often with a broad black lateral band, covering nearly the whole of the side of their bodies. Northern tamanduas are marked with a black V on their backs and a "vest" over their torso; southern tamanduas vary in appearance across their range, from a vested pattern like northern tamandua, to only partially vested, to entirely light- or dark-colored.
Behavior
Tamanduas are nocturnal, active at night and secreting away in hollow tree trunks and burrows abandoned by other animals during daylight hours. They can spend more than half of their time in the treetops, as much as 64%, where they forage for arboreal ants and termites. Tamanduas move rather awkwardly on the ground and are incapable of galloping like their relative, the giant anteater. Tamanduas walk on the sides of their clenched forefeet to avoid injuring their palms with their sharp claws.
Tamanduas manufacture a potent musk in their anal glands that they use for marking territory. They smear the strong-smelling secretions on rocks, trees, fallen logs, and other prominent landmarks to announce their presence to other tamanduas.
When threatened while in the trees, the tamandua will firmly grasp the branch with its hind limbs and tail and rear up to confront its attacker with slashing motions of its large, curved claws. When on the ground, it will protect its vulnerable hindquarters by backing against a tree or a rock and lashing out with its forearms.
The tamandua's small eyes afford limited vision. Instead of relying on their sense of sight, they primarily utilize their senses of smell and hearing to locate their insect prey. They use their sharp claws and powerful forearms to tear open the nest of a colony of termites and employ their elongated tongues, coated with sticky saliva, to extract the insects.
Female tamanduas reach sexual maturity at one year old and during their gestation period, which is around 160 days, typically have a single offspring. After being born, a tamandua pup rides on its mother's back and is left on a branch while its mother forages.
Conservation
The IUCN Red List considers both as species of concern. They are currently fairly common, yet face various threats. In Ecuador, tamanduas are killed as a precaution due to the erroneous belief that they attack domestic dogs. Elsewhere, they are hunted for meat or captured for the pet trade. They are also taken for the thick tendons in their tails, which are made into cordage.
Relationship with humans
The tamandua is frequently kept as an exotic pet. Tamanduas are also used occasionally by Amazonian Native Tribes as a form of biological pest control, utilizing them to rid their dwellings of termites and ants.
| Biology and health sciences | Xenarthra | Animals |
481843 | https://en.wikipedia.org/wiki/Uranium-235 | Uranium-235 | Uranium-235 ( or U-235) is an isotope of uranium making up about 0.72% of natural uranium. Unlike the predominant isotope uranium-238, it is fissile, i.e., it can sustain a nuclear chain reaction. It is the only fissile isotope that exists in nature as a primordial nuclide.
Uranium-235 has a half-life of 703.8 million years. It was discovered in 1935 by Arthur Jeffrey Dempster. Its fission cross section for slow thermal neutrons is about barns. For fast neutrons it is on the order of 1 barn.
Most neutron absorptions induce fission, though a minority (about 15%) result in the formation of uranium-236.
Fission properties
The fission of one atom of uranium-235 releases () inside the reactor. That corresponds to 19.54 TJ/mol, or 83.14 TJ/kg. Another 8.8 MeV escapes the reactor as anti-neutrinos. When nuclei are bombarded with neutrons, one of the many fission reactions that it can undergo is the following (shown in the adjacent image):
Heavy water reactors and some graphite moderated reactors can use natural uranium, but light water reactors must use low enriched uranium because of the higher neutron absorption of light water. Uranium enrichment removes some of the uranium-238 and increases the proportion of uranium-235. Highly enriched uranium (HEU), which contains an even greater proportion of uranium-235, is sometimes used in the reactors of nuclear submarines, research reactors and nuclear weapons.
If at least one neutron from uranium-235 fission strikes another nucleus and causes it to fission, then the chain reaction will continue. If the reaction continues to sustain itself, it is said to be critical, and the mass of 235U required to produce the critical condition is said to be a critical mass. A critical chain reaction can be achieved at low concentrations of 235U if the neutrons from fission are moderated to lower their speed, since the probability for fission with slow neutrons is greater. A fission chain reaction produces intermediate mass fragments which are highly radioactive and produce further energy by their radioactive decay. Some of them produce neutrons, called delayed neutrons, which contribute to the fission chain reaction. The power output of nuclear reactors is adjusted by the location of control rods containing elements that strongly absorb neutrons, e.g., boron, cadmium, or hafnium, in the reactor core. In nuclear bombs, the reaction is uncontrolled and the large amount of energy released creates a nuclear explosion.
Nuclear weapons
The Little Boy gun-type atomic bomb dropped on Hiroshima on August 6, 1945, was made of highly enriched uranium with a large tamper. The nominal spherical critical mass for an untampered 235U nuclear weapon is , which would form a sphere in diameter. The material must be 85% or more of 235U and is known as weapons grade uranium, though for a crude and inefficient weapon 20% enrichment is sufficient (called weapon(s)-usable). Even lower enrichment can be used, but this results in the required critical mass rapidly increasing. Use of a large tamper, implosion geometries, trigger tubes, polonium triggers, tritium enhancement, and neutron reflectors can enable a more compact, economical weapon using one-fourth or less of the nominal critical mass, though this would likely only be possible in a country that already had extensive experience in engineering nuclear weapons. Most modern nuclear weapon designs use plutonium-239 as the fissile component of the primary stage; however, HEU (highly enriched uranium, in this case uranium that is 20% or more 235U) is frequently used in the secondary stage as an ignitor for the fusion fuel.
Natural decay chain
Uses
Uranium-235 has many uses such as fuel for nuclear power plants and in nuclear weapons such as nuclear bombs. Some artificial satellites, such as the SNAP-10A and the RORSATs, were powered by nuclear reactors fueled with uranium-235.
| Physical sciences | Actinides | Chemistry |
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